README.md

metadata

license: creativeml-openrail-m
language: en
base_model: stable-diffusion-xl
tags:
  - stable-diffusion
  - text-to-image
  - generative
  - diffusion
metrics:
  - FID
pipeline_tag: text-to-image
library_name: diffusers

Model Card: JuggernautXL_v8Rundiffusion

Overview

JuggernautXL_v8Rundiffusion is a cutting-edge model built for stable diffusion-based generative tasks. This model is optimized for creating high-quality, detailed, and cinematic images using state-of-the-art diffusion techniques. It supports multiple LoRA (Low-Rank Adaptation) integrations and excels at generating diverse styles, including realistic, artistic, and abstract images.

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Features

• Model Type: Stable Diffusion XL (SDXL)
• Version: v8 Rundiffusion
• Key Capabilities:
  • Supports intricate artistic styles, cinematic imagery, and vibrant compositions.
  • Works seamlessly with additional LoRA files for fine-tuned styles.
  • Efficient memory usage and adaptability to various VRAM configurations.
• Applications:
  • Creative Design: Digital art, concept design, and visual effects.
  • Marketing: Ad campaigns, social media content creation.
  • Educational: Demonstrations of diffusion-based generative AI.

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Usage

This model requires a GPU with at least 16 GB VRAM for optimal performance.

Example Code:

from diffusers import StableDiffusionPipeline

# Load the model from Hugging Face
pipeline = StableDiffusionPipeline.from_pretrained("oieieio/juggernautXL_v8Rundiffusion")
pipeline.to("cuda")

# Generate an image
prompt = "a futuristic cityscape, cinematic lighting, ultra-detailed"
image = pipeline(prompt).images[0]

# Save the image
image.save("output.png")