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Wildfires in Northern Australia In early-October 2012, intense bushfires blazed in Australia’s Northern Territory in a remote area northeast of Elliott. The Moderate Resolution Imaging Spectroradiometer (MODIS) on NASA’s Aqua satellite captured this image of smoke from the fires streaming northeast on October 6, 2012. Red outlines indicate hot spots where MODIS detected unusually warm surface temperatures associated with fires. The light spots within the largest smoke plume are the tops of pyrocumulus clouds—tall, billowing clouds with well-defined borders. Such clouds form when intense heat from a fire pushes air rapidly upward, causing water vapor in the air to cool and condense into clouds. Pyrocumulus clouds sometimes generate intense thunderstorms that can either dampen the flames with drenching rains or propel the fire with strong winds. Wildfires are quite common in this part of Australia. Some areas, such as the higher rainfall savanna woodlands in the northern part of Northern Territory, burn nearly every year. - Wildfires burned in Australia during spring and summer 2012 - Earth Observatory. (2010, August 31). Russian firestorm: Finding a fire cloud from space. - Geosciences Australia. (2008, January). Understanding bushfire: trends in deliberate vegetation fires in Australia: Northern Territory. Accessed October 9, 2012. NASA image courtesy Jeff Schmaltz, LANCE MODIS Rapid Response Team at NASA GSFC. Caption by Adam Voiland.

http://www.nasa.gov/mission_pages/fires/main/world/Australia_amo_2012280.html

The NASA MESSENGER mission has found evidence for a significant amount of ice on Mercury, according to three papers published today in Science Express. Although Mercury is the closest planet to our Sun, MESSENGER data suggests that permanently shadowed pockets and craters near Mercury’s poles are cold enough to support water in the form of ice and other frozen volatiles. LASP developed and built the Mercury Atmospheric and Surface Composition Spectrometer (MASCS) instrument onboard MESSENGER. MASCS detects minerals on Mercury’s surface and determines the abundance and components of its thin atmosphere. MESSENGER instruments that contributed to the recent polar ice discovery include the Neutron Spectrometer and the Mercury Laser Altimeter. Since it entered Mercury’s orbit on March, 17, 2011, MESSENGER has collected multiple, independent lines of evidence for the existence of the planet’s icy, shadowy craters: It has used spectroscopy to measure hydrogen concentrations, which indicate the presence of water ice; measured the amount of light reflected by its polar deposits; and modeled the surface and near-surface temperatures of the polar regions. The resulting data suggest that ice is the primary component of Mercury’s north polar deposits, and that ice is exposed at the surface of some of the deposits and buried in others. For more information on the MESSENGER mission, please visit http://lasp.colorado.edu/messenger/.

http://lasp.colorado.edu/home/blog/2012/11/29/messenger-finds-polar-ice-on-mercury/

- Grades: 3–5 Protoceratops is the oldest of the horned dinosaur group, the suborder Ceratopsia in the order Ornithischia. Known from a great many skeletons collected in the Gobi Desert of Mongolia from late Cretaceous strata, about 90 million years old, it is one of the few dinosaurs whose very young stages are known. (The first dinosaur eggs to be discovered, in 1922, were once also attributed to Protoceratops but are now identified as those of Oviraptor.) Protoceratops, a herbivore, was about 2 m (6.5 ft) long and weighed 140 kg (300 lb) or more. It had a large turtle-beaked head, almost as long as the trunk of the body, and the flattened parietal and squamosal bones at the rear of the skull were flared out to form a crest, or frill. Unlike its descendants, it had no horns. Protoceratops has been found only in Asia, but all of its descendants are known only from North America. Bibliography: Dixon, Dougal, Dougal Dixon's Dinosaurs (1994); Glut, Donald F., Dinosaurs: The Encyclopedia (1995); Hatcher, John B., The Ceratopsia (1980).

http://www.scholastic.com/teachers/article/protoceratops

Story URL: http://news.medill.northwestern.edu/chicago/news.aspx?id=212698 Story Retrieval Date: 5/22/2013 11:37:59 AM CST The study of climates on hot, gaseous planets outside our solar system could help shine light on our own weather. Though climate circulation on "hot Jupiters" moves 90 degrees off from most weather patterns on Earth, scientists said this research could help create better climate models. "Hot Jupiters" are often many times the mass of the largest planet in our solar system - the planet they are named after. Gathering data about these planets is typically difficult because they are light years from our sun. But Northwestern University planetary scientist Nick Cowan helped develop a way to observe the 'hot Jupiters' from Earth by observing thermal radiation emitted by the planets. When paired with knowledge about the planet's orbit, scientists can determine climate patterns. Cowan and his team focus on "hot Jupiters" that have small orbits, meaning they are close to their suns and complete a full revolution in a matter of days. The composition of the planets and their vicinity to their stars mean the hemisphere closest to the sun heats relatively evenly. The side of the planet that faces away from the star is colder, so heat and weather patterns generally move on an east-west trajectory. This is different from the Earth. On Earth the poles are much colder than the equator, where our planet gets most of its radiation. So atmospheric circulation has strong north-south components, moving heat to the poles. But this is only one family of planets Cowan plans to observe through thermal radiation. He said he wants to gather large amounts of climate data from a variety of planet types, which could help create a universal climate model for planets. This could help us better understand our own climate on the Earth, he said.

http://news.medill.northwestern.edu/chicago/news.aspx?id=212698

Written by Stephen M. Edelson, Ph.D. Center for the Study of Autism, Salem, Oregon Stereotypy, or self-stimulatory behavior, refers body movements or repetitive movement of objects. This behavior is common in many individuals with developmental disabilities; however, it appears to be more common in autism. In fact, if a person with another developmental disability exhibits a form of self-stimulatory behavior, often the person is also labelled as having autistic characteristics. Stereotypy can involve any one or all senses. We have listed the five major senses and some examples of stereotypy. Sense Stereotypic Behaviors staring at lights, repetitive blinking, moving fingers in front of the eyes, hand-flapping tapping ears, snapping fingers, making vocal sounds rubbing the skin with one's hands or with another Vestibular (sense of balance) rocking front to back, rocking side-to-side placing body parts or objects in one's mouth, smelling objects, sniffing people. Why does stimming, or self-stimulation happen? Researchers have suggested various reasons for why a person may engage in stereotypic behaviors. One set of theories suggests that these behaviors provide the person with sensory stimulation (i.e., the person's sense is hyposensitive). Due to some dysfunctional system in the brain or periphery, the body craves stimulation; and thus, the person engages in these behaviors to excite or arouse the nervous system. One specific theory states that these behaviors release beta-endorphins in the body (endogeneous opiate-like substances) and provides the person with some form of internal pleasure. Another set of theories states that these behaviors are exhibited to calm a person (i.e., the person's sense is hypersensitive). That is, the environment is too stimulating and the person is in a state of sensory-overload. As a result, the individual engages in these behaviors to block-out the over-stimulating environment; and his/her attention becomes Researchers have also shown that stereotypic behaviors interfere with attention and learning. Interestingly, these behaviors are often effective positive reinforcers if a person is allowed to engage in these behaviors after completing a task. Intervention for stimming, or self-stimulation There are numerous ways to reduce or eliminate stereotypic behaviors, such as exercise as well as providing an individual with alternative, more socially-appropriate, forms of stimulation (e.g., chewing on a rubber tube rather than biting one's are also used to reduce these behaviors; however, it is not clear whether the drugs actually reduce the behaviors directly (e.g., providing internal arousal) or indirectly (e.g., slowing down one's overall motor movement). Copyright The purpose of this copyright is to protect your right to make free copies of this paper for your friends and colleagues, to prevent publishers from using it for commercial advantage, and to prevent ill-meaning people from altering the meaning of the document by changing or removing a few paragraphs. Reproduction kindly allowed by www.autism.org Visit their site for more useful resources. Click here for the full range of Asperger's and autism fact sheets at www.autism-help.org

http://autism-help.org/behavior-stimming-autism.htm

By Dr Ananya Mandal, MD Hearing loss of hearing impairment may be of two major types - conductive hearing loss and sensorineural hearing loss. A third type is a mixed type that has underlying symptoms of both these types of deafness of hearing loss. 1-6 Normal ear anatomy The normal ear consists of a narrow canal that lets in the sound waves. This is called the external ear or the ear canal. These waves enter the ear canal and strike the ear drum. The ear drum (called the tympanic membrane) is a membrane that vibrates as the sound waves hit it. These vibrations are passed to the three small bones (ossicles) inside the middle ear. These are called malleus, incus and stapes bones. The ossicles move to amplify the vibrations and pass them on to the inner ear. The inner ear contains a shell shaped organ called the cochlea. Within the cochlea are tiny hair cells all along the inner walls. These move in response to the vibrations and send a signal through the auditory nerve to the brain. Decibels hearing loss The normal hearing range is 0-20 decibels (dB). Around 30 dB are for whispers, 50 dB for average home noises and 60 dB for conversational speech. Sounds like jet engine noises are over 140 dB and are painful. Hearing loss is measured in decibels hearing loss (dB HL). - 25-39 dB HL means mild hearing loss (cannot hear whispers) - 40-69 dB HL means moderate (cannot hear conversational speech) - 70-94 dB HL is severe (cannot hear shouting) - more than 95 dB HL is profound (cannot hear sounds that would be painful for a hearing person) Types of hearing loss Types of hearing loss include conductive hearing loss, sensorineural hearing loss and mixed type. Conductive hearing loss In this the sound waves are unable to pass from the external ear into the inner ear resulting in a hearing loss. The most common reasons are due to: - blockage of the ear canal by ear wax - perforation of the ear drum - build-up of fluid due to an ear infection called glue ear Sensorineural hearing loss This occurs where the auditory nerve and other nerves that carry the information from the sounds heard to the brain are damaged due to age or injury. Hearing loss due to aging is called presbyacusis. After the age of 30 to 40, many people start to lose their hearing in tiny amounts. This increases with age and by 80 many people may have significant hearing impairment. Presbuacusis occurs when the sensitive hair cells inside the cochlea gradually become damaged or die. The initial symptoms include loss of high-frequency sounds, such as female or children’s voices and difficulty in hearing consonants, making hearing and understanding speech difficult. Ear injury is another common cause of hearing loss. This occurs due to damage caused by loud noises. The inner structures due to constant exposure to noise become damaged. Exposure to noise causes the hair cells inside the cochlea to be inflamed. Some drugs may also cause damage to the nerves of the ears leading to sensorineural hearing loss. These include notable antibiotics like aminoglycosides (Gentamicin, Amikacin etc.) Mixed type of hearing loss When people get both types together, the condition is termed mixed type of hearing loss. Reviewed by April Cashin-Garbutt, BA Hons (Cantab)

http://www.news-medical.net/health/Types-of-hearing-loss.aspx

total biodiversity of an area can be broken down into two hierarchical The number of functional types of organisms (e.g. carnivorous animals, Nitrogen (N)-fixing plants) or ecosystems (coniferous forest, prairie, tundra, marine intertidal). The number of functionally equivalent organisms/ genotypes within each functional type (e.g. the number of species of wood-rotting fungi). These organisms perform the same role in an ecosystem (e.g. moose and caribou are both large herbivores in boreal ecosystems; mussels and tunicates are both sessile marine filter feeders). (Huston, Ch.1 ) The basis for this division is that mechanisms that drive the diversity of functional types are different from the mechanisms that drive diversity among functionally equivalent organisms. For example, competition plays a significant role in determining diversity of functionally equivalent organisms but has little influence on functional diversity in an ecosystem. (Huston, Many other factors will influence the diversity of a system besides competition. These include evolutionary changes, geology, human history, environmental variability, disturbance and random population fluctuations. (Huston, These factors are further discussed in Part 3: Processes and Patterns have developed ways to characterize species diversity in a given diversity or alpha-diversity: refers to a group of organisms interacting and competing for the same resources or sharing the same environment. Measured as # of species within a given area. (Huston p.72; Whittaker, 1960, 1967; Fisher et al., 1943) or beta-diversity: refers to the response of organisms to spatial heterogeneity. High beta-diversity implies low similarity between species composition of different habitats. It is usually expressed in terms of similarity index between communities (or species turnover rate) between different habitats in same geographical area (often expressed as some kind of gradient). (Whittaker or gamma-diversity (Whittaker from figure 5.6 in Perlman, D.L. and Adelson, G., 1997. by permission of Blackwell Science, Inc. of species globally How many species are there on Earth? There is no definitive answer . Estimates fall between 1.5 and 30 million species of plants and animals. Another recent estimate claims that a more realistic number is 6 million. What we do know is that between 1.5 and 1.8 million species have been identified. The majority of species remain unidentified. Of the 34 known animal phyla, only one phylum lives exclusively on land while 33 are found in the ocean. Of those 33, 14 are found nowhere else on earth. Of the species that have been described, approximately: 000 of these are insects 41 000 are vertebrates 000 are plants The remaining species are comprised of invertebrates, fungi, algae and other microorganisms. The biological diversity of many ecosystems remains poorly explored, even today. These ecosystems include the deep ocean and the tree canopy and soil of tropical forests. more information on marine biodiversity, visit the Conservation

http://redpath-museum.mcgill.ca/Qbp/2.About%20Biodiversity/definition.htm

Slavery is an “abominable institution” that has plagued humanity throughout time. Slavery seems to be one of the few constants that we see throughout our shared past. The United States of America was built, to a great degree, on the backs of slaves. There has been a great deal of research into African slavery in America, from the buying of the slaves in Africa, to the voyage across the middle passage where so many people died, to the American markets where they were sold, to the living conditions that they endured in their new lives as slaves. Native American enslavement has not been nearly as researched or even known. Native Americans were kept as slaves in the new world. They were also slavers and slaveholders. Native Peoples were taken as war captives as well as kidnapped and taken from their homes. Some were sent to the West Indies, while others were kept as laborers and house servants. This research paper will focus on the question of whether or not Native slaves were treated differently from African slaves. The paper will also look at the concept of race in the seventeenth and eighteenth centuries to as to who was legally considered Indian or “Negro” by the courts and whether or not that made a difference in their daily lives. Adams, John. "Slavery and Race in Jeffersonian America." www.digitalhistory.uh.edu. www.digitalhistory.uh.edu/documents/documents_p2.cfm?doc=356 (accessed May 1, 2012). This source has writings from Adam’s promoting the slow abolition of slavery msainly for the safety of the country.. Hening, William. "Virginia Slave Laws." www.digitalhistory.edu. www.digitalhistory.uh.edu/documents_p2.cfm?doc=217 (accessed May 1, 2012). This source includes focuses on the confusion that was felt by whites about concepts of race. Jefferson, Thomas. "Native Americans and the American Revolutions." www.digitalhistory.uh.edu. www.digitalhistory.uh.edu/documents/documents_p2.cfm?doc=307 (accessed May 1, 2012). This source includes text from Jefferson about sending an expedition into Indian country to see if it was safe for settlers. Jefferson also writes about the brutality that he thinks will be necessary in Indian country to achieve the United States objectives in the region. Mittelberger, Gottlieb. "Immigration and Ethnic Diversity." www.digitalhistory.uh.edu. www.digitalhistory.uh.edu/documents/documents_p2.cfm?doc=227 (accessed May 1, 2012). Discusses the hardships endured by people on the voyage to America. "Objects in the Dark, 1636-1775 Tha Black Codes." www.hartford-hwp.com . www.hartford- hwp.com/HBHP/exhibit/02/2.html (accessed May 1, 2012). Includes a copy of the original “Black Codes” and information about what the documents meant within the context of the times. "The New Mafter and Miftress." www.pbs.org. May 1, 2012). Original document of a slave remembering her youth. "The Slave Experience:Living Conditions." www.pbs.org. www.pbs.org/wnet/slavery/experience/living/docs1.html (accessed May 1, 2012). A 1669 act that stated that masters would not be punished for killing the slaves that they "The Slave Experience:Living Conditions Original Documents." www.pbs.org. www.pbs.org/wnet/slavery/experience/living/docs8.html (accessed May 1, 2012). This source specifically focuses on the living condition of slaves in New England. Boston, Nicholas. "Slavery and the Making of America Living Conditions." www.pbs.org. www.pbs.org/wnet/slavery/experience/living/p_history.html (accessed May 1, 2012). This source discusses the living condition that slaves endured. Carocci, Max. “Written out of history: Contemporary Native American narratives of Enslavement.” Anthropology Today 25, no. 3 (June 2009): 18-22. Academic Search Complete, EBSCOhost (accessed April 3, 2012). This source discusses both Native slavery and how it has been “left out” of history. Gallay, Allan. Indian Slavery in Colonial America. The University of Nebraska Press, 2009. This book has several lengthy essays focused on the enslavement of Native Peoples from different perspectives. The first essay has the most information on Indian-Black issues in colonial Virginia. The essays on the Cherokee and Chickasaw are also helpful. Miles, Tiya. 2008. “The Narrative of Nancy, A Cherokee Woman.” Frontiers: A Journal of Women Studies 29, no. 2/3: 59-80. Academic Search Complete, EBSCOhost (accessed April 3, 2012). A discussion of one woman’s fight to be recognized as Indian and freed from slavery. Perdue, Theda. Slavery and the Evolution of Cherokee Society, 1540-1806. Knoxville: The University of Tennessee Press, 1979. (accessed May 1, 2012). This book focuses on Cherokees as slaves and slaveholders. This book also has a focus on Silver, Peter. Our Savage Neighbors. New York: W.W. Norton and Company, 2008. (accessed May 1, This book focuses on Indian-White relations in regards to the combination of religion and politics. Sturtevant(editor), Washburn(editor), William, Wilcomb. Handbook of North American Indians, 4 History of Indian-White Relations. Washington: Smithsonian Institution, 1998. (accessed May 1, 2012). This is a good starting point for finding an overview of Indian slavery and enslavement throughout the United States(has twelve pages on the subject).

http://www.austincc.edu/pgoines/johnston

Russia and the Soviet Union Under Lenin Russia and the Soviet Union Under Lenin The Communist regime that emerged from the civil war was characterized by chaos. Lenin and his associates had no experience of governing and had to create a system by trial and error. Lenin knew that he wanted an autocratic regime in which he would be sole dictator. He never considered establishing a parliamentary system; he believed this to be simply a rubber stamp for the capitalist forces of society. A major goal of the Russian revolutionaries was to incite similar revolutions throughout all of Europe; to destroy not just a type of government but an entire existing social and political order. To help bring this about, Lenin and his associates formed the International Communist Party, known as the Comintern, in 1919. The Comintern was characterized by rigid, uncompromising rules. Although the Socialist movement was strong throughout Europe, European Socialists were more moderate than Communists; they were on the whole satisfied with the greater degree of representation that ordinary citizens acquired during the nineteenth century. For example, by 1914 universal or near-universal adult male suffrage was the law in Britain, France, Germany, Italy, and elsewhere. Additionally, women acquired unprecedented freedom and political power during and immediately after World War I. In this accep- tance of the system, Lenin saw the defeat of everything he wanted to accom- plish; the Socialists of Europe simply were not prepared to go to the same extremes as the Communists. Through the Comintern, Lenin hoped to change this. Moscow controlled the Comintern from the early 1920s. In 1919, after World War I was over, the Great Powers met at Versailles to negotiate the peace. Russia took no part in the negotiations, but the country was nonetheless affected. The Russian territory that Germany had taken at Brest-Litovsk was made into independent nations; had Russia sent a delegation to Versailles, this might have been arranged differently. As matters stood, the Communist government refused to acknowledge the loss of the western territory until some time after Versailles. In the end, of course, Russia lost the fight to keep its land. War between Russia and Poland broke out in 1920. It did not last long. In March 1921 the peace treaty established the Russian-Polish border that would remain in place until 1939. Poland had become an enemy for Lenin to reckon with, for several reasons. First, it had a long-standing history of resentment toward Russian oppression. Second, it was a large nation with a large population, capable of holding its own in a struggle with Russia. Third, the Poles were fiercely anti-Bolshevik, in part because Poland was largely a Catholic nation and the Bolsheviks were atheists. The early 1920s in Russia can accurately be called “a Second Time of Troubles.” As a true Marxist, Lenin believed above all in policies that favored the workers. He also believed that industrialization was the key to Russia’s economic recovery. Therefore he instituted the New Economic Policy in 1921. It called for peasants to sell their surplus grain to the state at a fixed price in either money or kind (such as clothing or tools); the grain would be used to feed the urban industrial workers. The peasants reacted to the government orders in a way the Communists had not foreseen. Industry was crippled from the war and was not producing anything for the peasants to buy, so money was not useful to them; and the state rarely remembered to pay them in kind. Therefore, instead of working hard to provide the necessary surplus, they hoarded their grain, fearful of not having enough to feed their families. With no grain coming in from the country, the urban workers were going hungry; soon many of them were fleeing to the country in search of food. Severe droughts at this time led to widespread famine. Historians estimate that perhaps 6 million Russians died of starvation and disease during this period. In 1922, Russia was renamed to reflect the new government’s philosophy: it became the Union of Soviet Socialist Republics, called the Soviet Union or USSR for short. The twelve individual republics—including Georgia, Kazakhstan, and Russia—were equals, each with its own soviet, and all firmly under control of the dictator. The Communists made it clear to the old guard that there was no place for them in the new workers’ state. In the Soviet Union, the concept of private property disappeared. The wealthy were stripped of their homes, which were turned into apartment houses for workers, with the original owners perhaps being allowed to rent one room as their own family apartment. At least 2 million aristocrats packed what they could carry and fled to Western Europe. Those who stayed had to learn hard manual labor like all other Soviets. Add your own comment - Kindergarten Sight Words List - The Five Warning Signs of Asperger's Syndrome - What Makes a School Effective? - Child Development Theories - Why is Play Important? Social and Emotional Development, Physical Development, Creative Development - 10 Fun Activities for Children with Autism - Test Problems: Seven Reasons Why Standardized Tests Are Not Working - Bullying in Schools - A Teacher's Guide to Differentiating Instruction - Steps in the IEP Process

http://www.education.com/study-help/article/european-history-russian-revolution-soviet-unioin-lenin/

Floated items inside containers Date: 6 January 2004 Author: Russ Weakley A definition of Normal flow Normal flow is the way a document will display if you had no positioning or floating applied to elements. The content will flow down the page, starting with the first element in your document and finishing with the last element in your document. A definition of Float positioning A floated element is positioned within the normal flow, then taken out of the flow and shifted to the left or right as far as possible. More detailed information on floats can be found at Floatutorial Examples of floated element inside containers Below are four examples of floated items inside a container. The floated items are small red boxes. The containers are yellow boxes with dashed black borders. As mentioned above, floated elements (like the red box in the examples below) are taken out of normal flow, so their containers (the yellow containers) cannot determine their height. The yellow containers will stretch to fit content inside them – but only content that is in normal flow. The containers will ignore the height of the floated items. If the content inside the container is very short, the floated item may poke out the bottom of the container. If the content inside the container is long enough, the floated item may appear to sit correctly inside the container. If a container has no content inside, it has no height so it should not render the yellow background color at all. If there is only a floated item inside, then only this item will render. As you can see, the container will not render around the floated item. You can force the container to clear floated items inside. The container below has short content, but there is also a clearing element inside to push the container below the floated item. In this case, a <div> element is used, styled with “clear: both”. This is not an ideal solution as it involves adding additional elements. A simpler solution is to float the parent container as well. The container below has short content, but as the container is also floated, it wrappes around the red box. - Matt Brubeck: How to completely enclose a floated element in CSS2 - Tony Aslett: How To Clear Floats Without Structural Markup

http://www.maxdesign.com.au/articles/floatsample/

These classroom materials are written for students in grades 6-12 to deepen their understanding of the world and related social issues. These activities will raise awareness of individuals who make a difference and show how problem-solvers begin and carry out their work. Students are likely to have many questions after viewing The New Heroes series. These classroom resources will help them answer questions raised by these powerful stories. We hope to help students understand the following: - How do people really make a difference? - What is their approach? Their vision? - What are the steps or process to making a difference? - What kinds of people are good at this? Are there common characteristics and/or personal attributes that visionary problem-solvers have? - What do you care about? - What do the people around you need? - What is going on around you right now, in your own community? - What would you need to do now to get going? As a result of these learning experiences, students should be better able to: - Recognize both problems and potential solutions to local issues. - Recognize that there are different levels of entry into involvement in social issues. These resources contain two different types of materials to use in conjunction with the video series in your classroom. Lesson plans address themes that appear in multiple stories in the series and use short, specific video clips to explore the themes. Unit plans dig deeper into each of the four episodes in the series and use the full 20-minute video stories from the four episodes.

http://www.pbs.org/opb/thenewheroes/teachers/

Viral Hemorrhagic Fever (cont.) In this Article How are hemorrhagic fever viruses transmitted? Viruses causing hemorrhagic fever are initially transmitted to humans when the activities of infected reservoir hosts or vectors and humans overlap. The viruses carried in rodent reservoirs are transmitted when humans have contact with urine, fecal matter, saliva, or other body excretions from infected rodents. The viruses associated with arthropod vectors are spread most often when the vector mosquito or tick bites a human, or when a human crushes a tick. However, some of these vectors may spread virus to animals, livestock, for example. Humans then become infected when they care for or slaughter the animals. Some viruses that cause hemorrhagic fever can spread from one person to another, once an initial person has become infected. Ebola, Marburg, Lassa and Crimean-Congo hemorrhagic fever viruses are examples. This type of secondary transmission of the virus can occur directly, through close contact with infected people or their body fluids. It can also occur indirectly, through contact with objects contaminated with infected body fluids. For example, contaminated syringes and needles have played an important role in spreading infection in outbreaks of Ebola hemorrhagic fever and Lassa fever. What are the symptoms of viral hemorrhagic fever illnesses? Comment on this Specific signs and symptoms vary by the type of VHF, but initial signs and symptoms often include marked fever, fatigue, dizziness, muscle aches, loss of strength, and exhaustion. Patients with severe cases of VHF often show signs of bleeding under the skin, in internal organs, or from body orifices like the mouth, eyes, or ears. However, although they may bleed from many sites around the body, patients rarely die because of blood loss. Severely ill patient cases may also show shock, nervous system malfunction, coma, delirium, and seizures. Some types of VHF are associated with renal (kidney) failure. Reviewed on 11/22/2011 Viewers share their comments Viral Hemorrhagic Fever - Experience Question: Please describe your experience with viral hemorrhagic fevers. Viral Hemorrhagic Fever - Treatment Question: What was the treatment for your viral hemorrhagic fever? Viral Hemorrhagic Fever - Symptoms Question: What were the symptoms of your viral hemorrhagic fever? Viral Hemorrhagic Fever - Causes Question: What caused your viral hemorrhagic fevers? Get the latest health and medical information delivered direct to your inbox FREE!

http://www.medicinenet.com/viral_hemorrhagic_fever/page3.htm

What is chickenpox? Chickenpox is a highly contagious disease caused by the varicella-zoster virus (Herpes zoster). Who is at risk of chickenpox? Chickenpox occurs worldwide, affecting persons of all races, gender and age. However, it is largely a childhood disease with most cases occurring in children before 10 years of age. Once a person has had the chickenpox infection it is unlikely he or she will get it again as for most people one infection is thought to confer lifelong immunity. However, immunocompromised individuals are susceptible to the virus at all times and measures taken to either prevent or modify the course of the disease should be taken if there has been exposure to the virus. Although widespread chickenpox does not recur, the varicella virus remains in selected anterior horn cells of the spinal cord long term. It may be stimulated to reappear later as shingles, or herpes zoster infection. How do you get chickenpox? Chickenpox is highly contagious and is easily spread from person to person by breathing in airborne respiratory droplets from an infected person's coughing or sneezing or through direct contact with the fluid from the open sores. A person who is not immune to the virus has a 70-80% chance of being infected with the virus if exposed in the early stages of the disease. What are the signs and symptoms of chickenpox? In children, chickenpox usually begins as an itchy rash of red papules (small bumps) progressing to vesicles (blisters) on the stomach, back and face, and then spreading to other parts of the body. Blisters can also arise inside the mouth The spread pattern can vary from person to person. Also, depending on the individual case, there may be only a scattering of vesicles or the entire body may be covered with between 250 to 500 vesicles. The vesicles tend to be very itchy and uncomfortable. Some children may also experience additional symptoms such as high fever, headache, coldlike symptoms and vomiting and diarrhoea. Most adults who get chickenpox experience prodromal symptoms for up to 48 hours before breaking out in rash. These include fever, malaise, headache, loss of appetite and abdominal pain. The condition is usually more severe in adults and can be life-threatening in complicated cases. How is chickenpox diagnosed? Diagnosis of chickenpox is usually made on the presence of its characteristic rash (initial red papules that evolve into vesicles containing fluid) and that different stages of lesions are present simultaneously. A clue to the diagnosis is in knowing that the patient has been exposed to an infected contact within the 10-21 day incubation period. Patients may also have prodromal signs and symptoms. What is the treatment for chickenpox? For most healthy patients with chickenpox symptomatic therapy is usually all that is required. - Trimming children's fingernails to minimize scratching. - Paracetamol to reduce fever and pain (do not use aspirin in children as this is associated with Reye's syndrome). - Calamine lotion and/or oral antihistamines to relieve itching. - Consider oral aciclovir (antiviral agent) in people older than 12 years who may be at increased risk of severe varicella infections. Immunocompromised patients with chickenpox need intravenous treatment with the antiviral aciclovir. In cases of inadvertent exposure to the virus, varicella-zoster immune globulin if given within 96 hours of initial contact can reduce the severity of the disease though not prevent it. Chickenpox can now be prevented by vaccination with live attenuated varicella vaccine. Because the disease is usually uncomplicated and self-limiting in children, debate exists as to whether it should be given on a routine basis. The vaccine is currently not part of the immunisation programme for children in New Zealand. What are the complications from chickenpox? In healthy children, chickenpox infection is usually an uncomplicated, self-limiting disease. Problems that may arise in more complicated cases include: - Secondary bacterial infection of skin lesions caused from scratching - Dehydration from vomiting and diarrhoea - Exacerbation of asthma - Viral pneumonia Although the following complications may occur in healthy children with chickenpox, they are more commonly seen in immunocompromised and adult chickenpox cases. - Disseminated primary varicella infection; carries high morbidity - Central nervous system complications such as Reye's syndrome, Guillain-Barré syndrome and encephalitis - Thrombocytopaenia and purpura secondary to varicella infection Exposure to varicella virus may cause severe problems in pregnant women whom have not had chickenpox before. Chickenpox during pregnancy may cause viral pneumonia, premature labour and delivery and rarely maternal death. Also, approximately 25% of fetuses become infected. Offspring may remain asymptomatic, or develop herpes zoster at a young age without previous history of primary chickenpox infection. How to avoid spread of chickenpox A person with chickenpox is contagious 1-2 days before the rash appears and until all the blisters have formed scabs. This may take between 5-10 days. Children should stay away from school or childcare facilities throughout this contagious period. Adults with chickenpox who work amongst children, should also remain home. It can take from 10-21 days after contact with an infected person for someone to develop chickenpox. This is how long it takes for the virus to replicate and come out in the characteristic rash in the new host. Because of the serious complications that may occur in immunocompromised individuals and pregnant women, these people should avoid visiting friends or family when there is a known case of chickenpox. In cases of inadvertent contact, see your doctor who may prescribe special preventive treatment. Embedded external content may include advertising. On DermNet NZ: - Chickenpox– Medline Plus - Chickenpox – World Health Organization (WHO) - Chickenpox – Medscape Reference - Varicella – Immunisation Handbook 2002, Ministry of Health, New Zealand - Chickenpox (Varicella) – Immunisation Advisory Centre, University of Auckland - Chickenpox – emedicinehealth - Patient information: Chickenpox (The Basics) – UpToDate - Patient information: Chickenpox (Beyond the Basics) – UpToDate (for subscribers) Books about skin diseases: See the DermNet NZ bookstore

http://www.dermnetnz.org/viral/varicella.html

As you look up into a November sky right at nightfall, you may notice fewer bright stars than at other times of year. No, it’s not just the glare from Houston hiding most of the stars from view–there really are fewer bright stars in the November evening sky than in, say, February or August. To understand why, you need to understand the shape of our galaxy itself. Our galaxy, the Milky Way, is a barred spiral galaxy. Evidence indicates that the Milky Way, like many large galaxies, has a massive black hole at its center. A radio source designated Sagittarius A* could be the black hole itself. (The asterisk is part of the name, which is “Sagittarius-A-star”). Surrounding this black hole is a central bulge where older (and thus redder) stars predominate. The Bulge of our galaxy is not fully spherical but instead forms a bar a few thousand light years long. Branching out from this bulge are spiral arms which contain younger (bluer) stars and dust clouds out of which brand new stars form. Our solar system is about 26,000 light-years from the center to the edge, on the inside edge of the Orion Arm. The Orion Arm, in turn, is but a spur of the much longer Perseus Arm. The Milky Way is quite flat–over 100,000 light years wide but only 1,000 light years thick. The flatness of the galaxy means that most of its stars are near a certain plane in space. Of course, the galaxy is much thicker than our solar system, so we see our stellar neighbors suurounding us on all sides. The rest of the galaxy, extending off into the distance, appears to us as a hazy blur in the background, with individual stars (those fairly close to us) in the foreground. That hazy blur looked like spilled milk to the ancient Greeks, thus the name ‘Milky Way.’ We see more stars near that plane than far from it. What does this have to do with the dimness of a late November sky at dusk? Imagine observing our flat galaxy from our vantage point on Earth. When we face into the galactic plane, we see more bright stars, because there are more stars in that direction. When we face above or below that plane, we see fewer bright stars. Face west at dusk in late November and early December, and you’ll notice an enormous triangle of three bright stars, all bright enough to appear even in skies lit by Houston. These stars from the Summer Triangle, so called because it is up all night long from June through mid-August. This Triangle is also directly in the plane of the Milky Way. The constellation Sagittarius, which marks the center of the Galaxy, sets just after the Sun this time of year. Therefore, if you trace a path approximately from the point of sunset through the Summer Triangle, over to five stars in an ‘m’ shape in the North (that would be Cassiopeia, the Queen), and then over to the northeastern horizon. This is the plane of the Milky Way across late autumn skies at dusk . Turn to the south, and you face below the galactic plane (as we’ve arbitrarily defined ‘above’ and ‘below’). Here is a vast region of sky almost void of bright stars. One exception is Fomalhaut, low in the southeast at dusk tonight. Also, Houstonians with a very clear southern horizon can see Achernar very, very low in the south on December evenings. But that’s about it. There are many fewer bright stars in this direction than towards the Summer Triangle. By the way, the brilliant object in the east at dusk tonight, and high in the southeast as dusk in December, is Jupiter. It doesn’t count as a bright star for this sector of the sky. The Celestial Sea When ancient Mesopotamians looked up into the dim skies you see at dusk tonight, they imagined the Persian Gulf south of them extended up into the sky, forming a ‘Celestial Sea’. They therefore placed many water-themed constellations in this part of the sky. Zodiacal constellations here include Pisces, the Fish, and Aquarius, the Water-Carrier. Even Capricornus, the Goat, has the tail of a fish because he originally represented Ea, the ancient Babylonian god of the waters. Under Pisces is the sea monster Cetus, while Piscis Austrinus, the Southern Fish, drinks the water that Aquarius pours. Eridanus, the River, rises in the southeast, flowing from Orion’s foot into this vast ‘sea.’ Contrast this vast, dim region with the much brighter swath of stars that rises in the east later this evening (9-10 pm in late November, earlier in December). This region of sky includes the brilliant pattern Orion, the Hunter, as well as Sirius, the brightest star we ever see at night. When these stars rise, we are beginning to face back into the plane of our galaxy, this time looking into our own arm of galaxy at the stars right ‘behind’ the Sun. (This is why our arm of the Milky Way is called the Orion Arm.) Winter evening skies are much brighter than those of late autumn.

http://blog.hmns.org/tag/black-hole/

1. Outdoor Observations Take the class outside to identify where hummingbirds might be able to find nectar, insects, spiders, and tree sap. Encourage students to draw and write about the places they find in a field notebook. 2. Food Energy Research Challenge students to learn more about food energy by posing intriquing questions: How do the foods we eat become energy? Why is eating high-energy foods important? What food choices give you the energy you need each day? What kinds of food are high energy fuels? What food choices are "fizzle fuels"--foods that are a quick source of energy but that spurt of energy fizzles quickly? How much food do humans need to keep our engines running 24 hours a day, seven days a week, 365 days a year? links to explore: - Learn about how hummingbirds also help their partners, the sapsuckers! - Discover how brainy hummers are when it comes to remembering where to get the best food! - Observe the flowers that fuel migration and predict what they are a good "fit" for hummingbirds! 3. Take it Live: Track Hummingbirds' Spring Migration With Journey North As hummingbirds spread throughout their breeding range, continue to predict when and where they will travel.

http://www.learner.org/jnorth/tm/humm/sl/4/ts.html

How does right brain or left brain dominance affect how a student performs in the classroom? Most classroom teaching styles use left brain strategies. This tends to favor left brain dominant students, and can make learning difficult for right brain dominant students. Left brain students are good at linear and sequential processing, such as involved in language and math. Left brain students are also good at planning and following directions. These students easily learn information in lecture-style, teaching approach. Right brain students processes information more holistically. They learn by understanding the big-picture, not the details. They tend to be visual, not language oriented. This means they have more difficulty following a lecture-style teaching approach. Right brain students need to know why they are doing something. Right brain students can benefit from reviewing material before class to understand the bigger picture, and to understand the context for details that will be taught in class. Left brain students can easily express themselves in words. This is a large part of what is expected in class participation and in assignments. Right brain students may know what they want to say, but often have trouble finding the right words. A left brain student tends to be good with symbolic language and mathematics, and can easily memorize vocabulary words or math formulas. A right brain student needs to see, feel, or touch the real object. Right brain students prefer hands-on activities, and need to draw out a math or other problem to understand it. They also need diagrams or illustrations to help visualize the problem or solution. Right brain students learn visually, not by listening to a lecture-style class. They must take extensive notes, and use diagrams and drawings to make information more visual, to facilitate learning the information. They also need to make a mental images of things they hear or read in order to remember the information. Left brain students are good note takers and list makers. They are also good at planning and scheduling. This means they are good at completing assignments. Right brain students tend to approach things randomly. They tend to not make study schedules, and jump around from one task to another without regard to priorities. Right brain students may be late with an assignment, not because they weren't working hard, but because they were working on a lower priority assignment. Right brain students need extra effort in reading instructions to ensure they understand the assignment. They also extra effort in making assignment lists and study schedules. Left brain students are better at writing and spelling, since it involves sequencing and organizing of letters and words. Right brain students require more time to write a paper, and require more revisions to get it to say what they want to say. Right brain students must also rely more on spelling checkers and proof reading for their assignments. Right brain students tend to be more creative, but have more trouble than left brain students with the mechanics of writing and communicating.

http://kidport.com/RefLib/Science/HumanBody/NervousSystem/BrainHemispheres.htm

Computer programming has changed considerable since the days of FORTRAN and BASIC. Today, tools like Scratch can be used to engage and teach even young students fundamental programming concepts. As described on their About page: Scratch is a new programming language that makes it easy to create your own interactive stories, animations, games, music, and art — and share your creations on the web. Scratch is designed to help young people (ages 8 and up) develop 21st century learning skills. As they create Scratch projects, young people learn important mathematical and computational ideas, while also gaining a deeper understanding of the process of design. We have developed a six day Scratch unit for 7th and 8th grade math. One class period of each six-day cycle is devoted to Scratch, effectively spreading the lessons out over a six week period. During their course of study, students learn simple terminology, are introduced to the principles of object-oriented programming, and create original animations and games that are uploaded to our Scratch Web Gallery. The complete unit is available for download as a zip file and includes objectives, activities, materials, and homework. Alternatively, the major handouts can be found in this unit summary and the “lecture notes” have been combined into a Slideshare presentation: Scratch is much more powerful than most users realize and a six lesson unit, even when taught over an extended period of time, provides students with a good overview that leaves room for future exploration. Scratch enthusiasts should visit the Scratch website, Learn Scratch, and Collen Lewis’ Scratch page for additional project ideas, and encourage their students to learn a bit of programming; with Scratch, it really can be fun and easy! UPDATE: April 2010 Sergio González, a math teacher working at the Luis de Camoens school in Ceuta, Spain, recently updated and translated our six day algebra-geometry unit into Spanish, created video tutorials for some of the lessons, and assembled everything into a zipped package that can be imported into a LMS such as Moodle. The complete unit is now available for download in Spanish as a zip file and this open access sample class illustrates how Scratch could be presented in Moodle.

http://pwoessner.com/scratch-programming/

A new study by a team including scientists from the National Institute of Standards and Technology (NIST) indicates that thin polymer films can have different properties depending on the method by which they are made. The results suggest that deeper work is necessary to explore the best way of creating these films, which are used in applications ranging from high-tech mirrors to computer memory devices. Thin films spread atop a surface have many applications in industry. Inexpensive organic solar cells might be made of such films, to name one potential use. Typically they're made by dissolving the polymer, and then spreading a small amount of the liquid out on a surface, called a substrate. The solution becomes a film as the solvent dries and the remainder solidifies. But as this happens, stresses develop within the film that can affect its structure. Manufacturers would like to know more about how to control these stresses to ensure the film does what they want. But scientists who study film formation often use a different method of casting films than a manufacturer would. One method used in industry is "flow coating"—similar to spreading frosting across a cake. Another method is "spin casting"—placing a drop of liquid on a substrate that spins rapidly and spreads the droplet out evenly by centrifugal force. Both methods create smooth films generally, but the team decided to examine whether the two methods create different effects in finished films consisting of a self-assembling block copolymer. "It's an important question because some proposed applications intend to take advantage of these effects," Douglas says. The team's comparison led to results that surprised them. Although the rapid spinning of spin casting is very dynamic, suggesting it would convey more stress to the resulting film, it actually led to fewer residual stresses than flow coating did. As previous studies have shown that leftover solvent can lead to stresses in the film, the team's new theory is that because the solvent evaporates from the developing film more slowly in flow coating, this solvent discourages the film solids from arranging themselves into the equilibrium structure. For one example, the practical benefits of this understanding could help manufacturers who propose making computer memory devices from thin films in which the solids arrange themselves as tiny cylinders in the film. Such devices would require the cylinders to stand on end, not lay down flat. "We find we can get them to stand up much more easily with one casting method than another," Douglas says. "If we can get better results simply by varying the mode of film casting, we need to explore more deeply what happens when you make films by different methods." Explore further: Two collider research teams find evidence of new particle Zc(3900) More information: Soft Matter, Mar. 21, 2012. doi:10.1039/C2SM07308K.

http://phys.org/news/2012-05-fabrication-method-affect-block-copolymer.html

I've always wondered how exactly the missing value ( R is represented under the hood. Last weekend I was working on a little project that gave me enough excuse to spend some time on finding this out. So, I descended into the catacombs of R and came back with some treasure. In short: - A missing integer is repesented by the largest negative (4 word) signed integer that can be represented by your computer. - A missing double(real number) is represented by a special version of the default NaN(Not a Number) of the IEEE standard. A special role is given to the number 1954 here, but why? Read on if you want to dig a little deeper. As you may know, a lot of R's core is written in the C language. However, an int variable in C does not support the concept of a missing value. So, what happens in R is that a single value of the integer range is pointed out as representing a missing value. In this case it is C macro from limits.h) which determines the largest negative value that can be represented by a int variable in C. On most computers, an int variable will be 32 bits (4 8-bit words). To make things easier, we'll assume that's always the case here. Since 1 bit is reserved for the sign, the range of representable numbers is (The range is asymmetric because 0 occupies the place of one positive number). Now let's compare this with R's integer range. The maximum integer is easily found, since it is stored in the hidden > .Machine$integer.max 2147483647 So this corresponds with INT_MAX. The largest negative integer is not present in .Machine but we can do some tests: 1 2 3 4 5 6 7 8 9 10 11 12 13 14 15 16 17 18 19 # store the one-but-least C-integer. The L in the end forces the number # to be "integer", not "numeric" x <- -2147483647L typeof(x) "integer" # adding an integer works fine, since we move further into the range: typeof(x+1L) "integer" # substracting an integer gives a warning telling us that the result is out-of-range: > typeof(x-1L) "integer" Warning message: In x - 1L : NAs produced by integer overflow # substracting a non-integer 1 ("numeric") yields a non-integer: > typeof(x-1) "double" The result is out of R's integer range. The integer range of R is : one integer less than you get in C. So by sacrificing only one of your four billion two hundred ninety-four million nine hundred sixty-seven thousand two hundred ninety-five integers, you get the truly awesome feature of computing with missing values. To explain how real () missing values are represented, we first need to spend a few words on the double type. A double is short for double precision and it is the variable type used to represent (approximations to) the real numbers in a computer. Basically, a double represents a rounded real number in the following notation (see also the wikipedia article): The sign is represented by 1 bit, the exponent by 11 bits and the mantissa by 52 bits, so we have 64 bits in total. The special value NaN (and also Inf) is coded using values of that are not used to represent numbers. NaN is represented by 0x7ff (hexadecimal) and . The important thing is that it does not matter what the value of is when representing NaN. This leaves developers with lots of room in the mantissa to give different meanings to R the developers chose in the mantissa to represent C-level function called R_IsNA detects the 1954 in A funny question is why did the R developers choose 1954? Any ol' number would have been fine. Was it because - It's the year of birth of one of the developers? (I couldn't find a match here) - Alan Turing died in 1954? (macabre) - Because president Eisenhower met with aliens in 1954? (ehm...) - In 1954 Queen Elisabeth II became the reigning monarch of Australia? (well...) Leave an answer in the comments if you have a better idea...

http://www.markvanderloo.eu/yaRb/2012/07/08/representation-of-numerical-nas-in-r-and-the-1954-enigma/

Video 1. Statistics As Problem Solving Consider statistics as a problem-solving process and examine its four components: asking questions, collecting appropriate data, analyzing the data, and interpreting the results. This session investigates the nature of data and its potential sources of variation. Variables, bias, and random sampling are introduced. Go to this unit. Video 2. Data Organization and Representation Explore different ways of representing, analyzing, and interpreting data, including line plots, frequency tables, cumulative and relative frequency tables, and bar graphs. Learn how to use intervals to describe variation in data. Learn how to determine and understand the median. Go to this unit. Video 3. Describing Distributions Continue learning about organizing and grouping data in different graphs and tables. Learn how to analyze and interpret variation in data by using stem and leaf plots and histograms. Learn about relative and cumulative frequency. Go to this unit. Video 4. The Five-Number Summary Investigate various approaches for summarizing variation in data, and learn how dividing data into groups can help provide other types of answers to statistical questions. Understand numerical and graphic representations of the minimum, the maximum, the median, and quartiles. Learn how to create a box plot. Go to this unit. Video 5. Variation About the Mean Explore the concept of the mean and how variation in data can be described relative to the mean. Concepts include fair and unfair allocations, and how to measure variation about the mean. Go to this unit. Video 6. Designing Experiments Examine how to collect and compare data from observational and experimental studies, and learn how to set up your own experimental studies. Go to this unit. Video 7. Bivariate Data and Analysis Analyze bivariate data and understand the concepts of association and co-variation between two quantitative variables. Explore scatter plots, the least squares line, and modeling linear relationships. Go to this unit. Video 8. Probability Investigate some basic concepts of probability and the relationship between statistics and probability. Learn about random events, games of chance, mathematical and experimental probability, tree diagrams, and the binomial probability model. Go to this unit. Video 9. Random Sampling and Estimation Learn how to select a random sample and use it to estimate characteristics of an entire population. Learn how to describe variation in estimates, and the effect of sample size on an estimate's accuracy. Go to this unit. Video 10. Classroom Case Studies, Grades K-2 Explore how the concepts developed in this course can be applied through a case study of a K-2 teacher, Ellen Sabanosh, a former course participant who has adapted her new knowledge to her classroom. Go to this unit. Video 11. Classroom Case Studies, Grades 3-5 and 6-8 Explore how the concepts developed in this course can be applied through case studies of a grade 3-5 teacher, Suzanne L'Esperance, and a grade 6-8 teacher, Paul Snowden, both former courses participant who have adapted their new knowledge to their classrooms. Go to this unit.

http://www.learner.org/resources/series158.html

Lexile Framework for Reading Lexile measures rank the reading level of a student and the difficulty of a text on a single scale and are correlated to scores on Standards of Learning (SOL) reading assessments in grades 3-8. Teachers and parents use Lexile measures to select books that are likely to improve the reading and comprehension skills of students. Lexile measures are reported with the scores of students taking SOL reading tests in grades 3-8. Virginia’s Focus on Literacy Achievement and the Lexile Framework for Reading Students are more likely to read books that match their current reading level. By using Lexile measures, teachers can assign and recommend books and other reading materials that will help students develop stronger reading skills. Parents can use Lexiles to select texts that reinforce what teachers are trying to accomplish in the classroom. It is important to note that the Lexile measure does not address the content or quality of the book. Many other factors affect the relationship between a reader and a book, including its content, the age and interests of the reader, and the design of the actual book. The Lexile measure is a good starting point in the book-selection process, but parents and educators should always consider these other factors when making a decision about which book to choose. The Virginia Lexile Framework Map (PDF) provides a sampling of titles with corresponding Lexile levels. The framework illustrates how Lexiles can be used to support students and to monitor their progress as they learn to become better readers. Additional book titles have been measured on the Lexile scale and can be found under Look up a Book. Books listed are not endorsed or recommended by the Virginia Department of Education. These resources and more may be accessed by using the menu to the right. The menu includes links to detailed information for parents and educators, frequently asked questions and related assessment and school improvement resources.

http://www.doe.virginia.gov/testing/scoring/lexile/index.shtml

Before Federation: To 1900 Rum, Shinplasters, Holey Dollars & More When the colony of New South Wales was established, no provision was made for an internal currency. There were no banks for some time and of course no central bank. This produced a period of chaotic currency arrangements and experiments in the colony. In the early years, the colonists had to rely on barter – the exchange of produce, goods and services – and other makeshift currencies such as rum (as all spirits were then called). Coin was very scarce in the early years with England also suffering from a coin shortage. The Spanish dollar was a major international currency of the time. A shipment of Spanish dollars was sent from England in 1792. Coins from a range of other countries were also used such as Dutch guilders and ducats, Indian mohurs and rupees and Portuguese johannas. But much of this coin left the colony as a result of trade with visiting merchant ships. Governor Macquarie also tried to remedy the coin shortage by the creation of the holey dollar. He had the centres of Spanish dollars punched out, leaving a ring (the holey dollar) valued at 5 shillings and a dump or core valued at one shilling and threepence. Despite stiff penalties for exporting this coin, it remained scarce. The Commissariat, which controlled the issue of stores to troops and convicts, bought goods produced by private enterprise, paying for them with store receipts. These receipts served as a medium of exchange up to the 1820s but they were for unwieldy amounts. The Commissariat began to issue its own notes, which were equivalent to those of the Bank of New South Wales, established in 1817. Because of the persistent shortage of coin and the limitations of other currencies, promissory notes or IOUs also soon came into general use. Squatters' cheques were a form of IOU in specific districts. These IOUs, however, were vulnerable to forgery and had no collateral backing. Promissory notes became known as shinplasters, a term describing a paper currency thought to be only worth soaking in vinegar as a poultice for bruises. These notes mostly fell apart quickly in the pockets or boots of customers, thereby saving the issuer from having to redeem them in coin. From the late 1840s, copper tokens were issued by businesses as small change to relieve the coin shortage but, like paper IOUs, had no backing or official guarantee. With no adequate solution to the currency problem, the British Government had legislated a sterling currency for the colony in 1825. 'Pounds, shillings and pence' remained in place as the basis of Australian currency until the introduction of the decimal system in 1966.

http://www.rba.gov.au/Museum/Displays/1788_1900_before_federation/currency_chaos.html

In this section you will learn how to develop your first Java program and then compile and test the same on your development machine. Write first Program in Java In Java, all source code is written in plain text files .java extension. These source files are then .class files by the .class file does not contain code that is native to your processor; instead it contains bytecodes that is the machine language for the Java Virtual Machine. You can write a Hello World program as your first Java program. Hello world program is the first step of java programming language. To write the Hello world program you need simple text editor like note pad and make sure the JDK must be install on your machine for compiling and running that program. Learn to write Hello World Java program at http://www.roseindia.net/java/master-java/helloworld.shtml Compile and Test the Program Java provides javac compiler to compile the Java program. To compile your HelloWorld program, open the command prompt. Go to the directory where you have saved the file and issue the following command to compile the program: Recommend the tutorial

http://www.roseindia.net/java/gettingstartedwithjava/javafirstprogram.shtml

Powdery mildew fungi are ideal for use in a laboratory exercise on fungal classification, fungal spore types, host-parasite interaction during an ecology unit, or a study of host range. Students will discover how many different plant hosts they can find that are infected by the same genus of a powdery mildew fungus, or how many different genera of powdery mildew fungi can be found on the same plant host. This exercise demonstrates the diversity that exists within a fungal order. With a good collection of leaves infected with different powdery mildew fungi, students learn to use a written key and/or an illustrated key (or could even make their own key) to identify the powdery mildew fungus to genus. Since powdery mildew fungi reproduce by means of two spore types: asexual spores (conidia) and sexual spores (ascospores), discussions of the types of reproduction in fungi would be facilitated. Powdery mildews are one of the most common, conspicuous, widespread, and easily recognizable plant diseases. As a group, powdery mildew fungi infect many species of plants, including many trees and shrubs, numerous ornamentals, vegetables, cereals, grasses, and even weeds. However, individual species of powdery mildew fungi are usually very host specific. For example, the species of fungus causing powdery mildew on pumpkin is different from that causing the disease on roses. Extensive losses in plant growth and crop yield occur annually due to powdery mildew. The primary sign of powdery mildew is grayish white, powdery blotches on leaves. Usually, powdery fungal growth appears first on the upper leaf surface (FIGURE 1). Eventually the entire leaf may become covered with mildew (FIGURE 2). Powdery mildew fungi also can infect other parts of the plant and may cause distortion and stunting of shoots, leaves (FIGURE 3), and flowers (FIGURE 4), and russeting on fruit (FIGURE 5). The type and extent of symptoms vary depending on the combination of powdery mildew fungal species and host plant species involved. Late in summer and fall, the sexual stage of many species of powdery mildew fungi, the cleistothecia, is visible as black or brown, pinhead-sized, spherical specks among the white to grayish mildew mycelium in the older infected areas on the leaves of many plants (FIGURE 6). (These are the leaves that should be collected and pressed flat between newspaper for storage and use in this laboratory exercise.) The powdery mildew diseases of various crops and other plants are caused by many different species of fungi grouped into six main genera in the order Erysiphales, an order that includes a single family, the Erysiphaceae. The fungi causing powdery mildews are obligate parasites (or biotrophs), meaning that they cannot be cultured on nutrient media and require a living host to grow. These fungi reproduce by means of two spore types: asexual spores called conidia and sexual spores called ascospores. The conidia are usually barrel-shaped or oval and are usually formed in chains (FIGURE 7) at the ends of specialized hyphae called conidiophores produced from the mycelium growing on the surface of a plant’s leaves, stems, flowers or buds. The combination of the mycelium, conidia and conidiophores gives the leaf surface a powdery appearance from which the name powdery mildew is derived. The fungi are spread when the conidia are carried by air currents to new plant surfaces. Under favorable conditions, the fungal spores germinate (FIGURE 8). Unlike most fungal parasites that invade plant tissues, most powdery mildew fungi grow superficially on top of the leaf surface (FIGURE 9). Fine, thread-like infection pegs penetrate the epidermal cells of the leaves and form haustoria (FIGURE 10). Haustoria are specialized hyphae for nutrient absorption from the plant cells. The plant is damaged by the loss of nutrients to the fungus, disruption of photosynthesis, and premature death of leaves or other infected plant tissues. When environmental conditions or nutrition become unfavorable for growth (usually later in the growing season), the fungus shifts to the sexual stage and produces cleistothecia (sing. cleistothecium). The cleistothecia are closed, thick-walled, tiny, black, spherical structures (white to tan when young) that house sacs called asci (sing. ascus) (FIGURE 11 and FIGURE 12). The oval sexual spores, called ascospores, are produced within the asci. Fungi that produce ascospores in asci belong to the Ascomycetes, a fungal group that includes yeast, morels, and many plant parasites. Cleistothecia have ‘arm-like’ appendages that radiate out from their outer surface. Inside each cleistothecium is a single ascus or many asci. Currently, the powdery mildew fungi are classified to genus based on the number of asci contained in the cleistothecium and on the morphology (physical appearance) of the hyphal appendages (arms) growing out of the wall of the cleistothecium (FIGURE 11, FIGURE 12, and Key to Genera of Powdery Mildew). Recently, Saenz and Taylor (1999) have proposed that powdery mildew genera be reclassified according to the phylogeny (an estimate of evolutionary relationships) inferred from the DNA sequence of the ribosomal ITS region and a number of morphological characteristics. If this occurs, the morphology of conidia and conidiophores will be used as the primary characters for classification

http://www.apsnet.org/edcenter/K-12/TeachersGuide/PowderyMildew/Pages/Background.aspx

Dehydration is the excessive loss of body fluids. Our bodies need a certain amount of fluid daily to maintain all the body systems with the minimum considered to be 4 glasses of fluid per day (the ideal amount is considered to be 8 glasses). However, basic requirements will vary with your age and the amount of activity you undertake. Very active people may need 2 to 3 times the normal limit. If we lose more fluid than we replace the result is dehydration. Normally, the fluid depletion is signalled by thirst and is replenished by drinking. If the fluid lost is not replaced within a few days, severe dehydration develops, and can lead to changes in the body’s chemistry, kidney failure and may even be fatal. This severe form of dehydration is considered a medical emergency. There are also a number of conditions, types of surgery and medications that can cause dehydration. Your doctor can determine the cause of dehydration and treat it, and may admit you to hospital for the administration of intravenous fluids if dehydration is severe. Dehydration can be prevented by drinking plenty of fluids if you are in a hot climate, participating in strenuous exercise or suffering from diarrhoea. If you are vomiting uncontrollably, try to sip small amounts of fluid in between bouts. Last Reviewed: 11 July 2001

http://www.mydr.com.au/nutrition-weight/dehydration

"Mountains are [therefore] amazing places to observe vegetation changes in response to climate warming," Lenoir added. The new study drew on nearly 8,000 historical surveys of the mountains in and around France, some stretching back to 1905. Temperatures in these mountains—which include the Western Alps, Pyrenees, and Massif Central—crossed a threshold around 1985, the researchers say. Before that year the region shows no clear trend in climate changes. But since then the mountains have been warming, and plants began moving in sync with rising temperatures. The scientists looked at the movements of 171 species in forests on the lower slopes, from sea level up to 8,500 feet (2,600 meters). While earlier studies had focused on plants in high altitudes that are known to be more sensitive to temperature changes, the new work found that even common plants at lower elevations are feeling the heat. The team also discovered that different types of plants are moving at different rates. "Long-lived plants like trees or shrubs did not show a significant shift, whereas short-lived species like herbs showed a strong upward shift in elevation," Lenoir said. "This may imply profound changes in the composition and the structure of plant communities and on the animal species they interact with," he added. "It may disrupt ecosystems." Out of Room Steven Running is an ecologist at the University of Montana in Missoula who was not involved in the new study. "Theory would predict that warming temperatures will allow plant distributions to expand into cooler, higher mountain elevations," he said. "This paper confirms the theory. With a large population analysis, [it is] the best evidence published so far that plant distributions are rising [in altitude]." The shifts in elevations of some plants and not others "shows very much that the ecosystems are already evolving away from the ecosystems as we know them," added Cynthia Rosenzweig of the NASA Goddard Institute for Space Studies in New York City. With shorter-lived species, "we're seeing this rapid response" to warming, Rosenzweig said. While it shows these plants are adapting to the changes, she noted, "they're going to have to keep moving up and up—and eventually run out of room." And the longer-lived species that aren't moving may also be headed for trouble, she said. "There's concern that they aren't adapting and may not have spread their seeds far enough" to shift their ranges to cope with continued warming. SOURCES AND RELATED WEB SITES

http://news.nationalgeographic.com/news/2008/06/080626-plants-warming_2.html

Discover the cosmos! Each day a different image or photograph of our fascinating universe is featured, along with a brief explanation written by a professional astronomer. November 12, 1997 Explanation: El Niño is a temporary global climate change resulting from unusually warm water in the central Pacific Ocean. El Niño can cause unusual or severe weather for some locations over the next few months. Warm water is shown in white in the above false-color picture taken by the orbiting TOPEX/Poseidon satellite in late October. The Pacific Ocean is color coded by sea surface height relative to normal ocean levels. The large white area represents a mass of warm water 30 times greater than all the Great Lakes, flowing toward the Americas. Although El Niños occur every decade or so, this year's is the first ever predicted. The cause and full effects of El Niños are still under study. Authors & editors: NASA Technical Rep.: Jay Norris. Specific rights apply. A service of: LHEA at NASA/ GSFC &: Michigan Tech. U.

http://apod.nasa.gov/apod/ap971112.html

Earthquakes shake the ground because fault rupture releases vibrations that radiate in the form of seismic energy. These earthquake waves, also called body waves, come in two distinct forms: Primary or "P" waves and Secondary or "S" waves. When body waves reach the free surface of the earth some of their energy is converted into complex surface waves that are trapped near the surface of the earth and produce generally lower frequency ground motions. P waves are compressive waves that do not produce much damage. They can move through any type of material and travel at almost twice the speed of S waves. High frequency P waves due not weaken or attenuate as rapidly as S waves so they retain higher frequencies when they arrive at seismic stations. In air, P waves take the form of sound waves and therefor move at the speed of sound, 330m/s at sea level. Some people even report hearing an earthquake (due to the higher frequency P waves vibrating or rustling objects) before they feel the S waves arrival. Typical speeds in Earth are faster: 1450m/s in water and 5000m/s in granite. S waves are shear waves that deform the ground perpendicular to their direction of travel. Unlike P waves, S waves are unable to pass through air and liquids such as water and magma. This is how we know Earth’s outer core is molten – these waves cannot pass through it! When S waves deform the ground, it causes lateral or shear (back and forth) forces on structures. Older buildings were constructed primarily to withstand gravity (vertical forces); therefore they are more prone to fail due to the strong lateral loading experienced during a big earthquake. - Both P and S waves are generated across a broad spectrum of frequencies. The higher the frequency, the faster the energy from the earthquake attenuates, or dissipates, with distance. Also, due to attenuation and geometrical spreading, locations close to the source of the rupture that caused the earthquake will receive more energy (and shaking) than more distant locations. - When P and S waves arrive at the surface of the Earth, part of the energy is trapped and guided by Earth's surface. Their behavior is different than for body waves, so we say they are converted to Surface Waves. The two types are Love waves and Rayleigh waves. Rayleigh waves produce a long rolling motion along the earth's surface much like the motion in a boat on the open sea. They travel a little slower than Love waves.

http://www.pnsn.org/outreach/about-earthquakes/eq-waves

What is Inflation? Inflation is the term used to describe a rise of average prices through the economy. It means that money is losing its value. The underlying cause is usually that too much money is available to purchase too few goods and services, or that demand in the economy is outpacing supply. In general, this situation occurs when an economy is so buoyant that there are widespread shortages of labour and materials. People can charge higher prices for the same goods or services. Inflation can also be caused by a rise in the prices of imported commodities, such as oil. However, this sort of inflation is usually transient, and less crucial than the structural inflation caused by an over-supply of money. Inflation can be very damaging for a number of reasons. First, people may be left worse off if prices rise faster than their incomes. Second, inflation can reduce the value of an investment if the returns prove insufficient to compensate them for inflation. Third, since bouts of inflation often go hand in hand with an overheated economy, they can accentuate boom-bust cycles in the economy. Sustained inflation also has longer-term effects. If money is losing its value, businesses and investors are less likely to make long-term contracts. This discourages long-term investment in the nation’s productive capacity. The flip-side of inflation is deflation. This occurs when average prices are falling, and can also result in various economic effects. For example, people will put off spending if they expect prices to fall. Sustained deflation can cause a rapid economic slow-down. The Reserve Bank is as concerned about deflation as it is about inflation. In New Zealand, however, it has historically been more usual for prices to rise. As Figure 1 shows, there have been only brief periods of deflation in the past 150-odd years, and these have been associated with economic depressions. The graph also shows that, once the economy had become established, New Zealand did not have sustained high inflation until the 1970s and 1980s. In the late 1980s the government gave the Reserve Bank responsibility for keeping inflation low and more stable than it had been. Statutory authority was provided in the Reserve Bank of New Zealand Act 1989, and the specifics were set out in a written agreement between the Governor of the Reserve Bank and the Minister of Finance. This ‘Policy Targets Agreement’ initially called for a reduction of inflation to 0–2 percent increase in the Consumers Price Index (CPI) by 1992. It has been revised several times since, and the current agreement, signed in May 2007, calls for inflation to be kept within 1 to 3 percent a year, on average over the medium term. This means that, as the graph shows, inflation can exceed the 1–3 percent target range in the short term. However, in the medium term it remains within that band, on average, and this means that the very high inflation rates of the 1960s and 1970s – which at times exceeded 18 percent per annum – do not occur. The effect of this arrangement is clear from Graph 2, in which inflation has remained within a narrow band. The Bank controls inflation through an economic tool known as the Official Cash Rate, covered in a separate sheet. There are various ways of measuring inflation. The one used in the Policy Targets Agreement is the CPI published by Statistics New Zealand. This records the change in the price of a weighted ‘basket’ of goods and services purchased by an ‘average’ New Zealand household. Statistics New Zealand weights and indexes the various items in the basket and forms the ‘all-groups’ index. The percentage change of this index is typically referred to as ‘CPI inflation’, and is usually expressed over both a quarterly and annual period. The contents of the basket are defined by Statistics New Zealand, which periodically reviews and re-weights them, using data obtained from their annual Household Economic Survey. This is necessary because the basket of goods and services purchased by the average household will change over time. The Reserve Bank has published an interactive inflation calculator on its website, at:http://www.rbnz.govt.nz/statistics/0135595.html This calculator allows users to input a sum of money and compare its value, in terms of the CPI and other measures for pre-CPI years, between any two quarters from 1862, to the latest quarter for which CPI figures are available.

http://www.rbnz.govt.nz/monpol/about/0053316.html

American educators have struggled for more than 40 years to define giftedness. Yet even now, there is no universally agreed upon definition of what it means to be gifted. U.S. federal law defines gifted students as those who perform or who show promise of performing at high levels in any one of five categories: general intellectual ability, specific academic aptitude, creative or productive thinking, leadership ability or visual/performing arts. Beyond that definition, there are no specific national criteria for identifying gifted and talented students nor does federal law provide funding or mandates for identification of these students or programming for them. This definition is left to the states. The result has been a wide variety of state definitions and methods for the identification of gifted children. Some states have specific definitions for giftedness, while others have none. Some states require programs for gifted students, while others do not. In other words, the availability of programs and services for gifted students depends for the most part on where a student lives and what state, school district or school he or she is in. There is debate over how to identify and measure giftedness, whether giftedness is innate (nature) or developed (nurture) and whether giftedness is driven by intelligence test results or through other indicators. These varying perspectives have led to much misinformation about gifted students and what programs for gifted students should look like. Here are 10 of the most common myths about gifted students and programs for the gifted: Myth No. 1: Intelligence is inherited and does not change. Gifted students, therefore, do not need any special services. All of us do inherit certain traits, intelligences and talents. But these need to be developed and nurtured throughout life for them to grow and reach their full potential. A beautiful flower inherits certain traits. But if it is not watered and fed and if it does not get the right amount of sunlight, it does not develop as it could. The same is true for gifted children. Myth No. 2: Giftedness can easily be measured by intelligence tests and tests of achievement. Giftedness is difficult to measure. This is why schools and school districts try so many different ways to identify gifted students. Tests are often culturally biased and may reflect ethnicity, socioeconomic status, exposure and experiences rather than true giftedness. Other children may be gifted but are not good at taking tests. They may not score well on standardized tests but may be gifted, especially in creative and productive thinking.

http://www.wyff4.com/news/national/My-view-Ten-myths-about-gifted-students-and-programs-for-gifted/-/9324256/17401522/-/4gf817z/-/index.html

Science Fair Project Encyclopedia Shintō (Japanese: 神道) is the native religion of Japan. It involves the worship of kami, which could be translated to mean gods, nature spirits, or just spiritual presences. Some kami are very local and can be regarded as the spirit or genius of a particular place, but others represent major natural objects and processes, for example, Amaterasu, the Sun goddess. The word Shinto was created by combining two Chinese characters (神道, read shen dao in modern Chinese). The first character means "divine" or "God", and can also be read as "kami" in Japanese. The second character means "way" or "path," and is the character used for the word "Taoism." Thus, Shinto literally means "the way of the kami." After World War II, Shinto lost its status of national religion; most Shinto practices and teachings, once given a great deal of prominence during the war, are no longer taught nor practiced today, and some remain largely as everyday activities without religous connotations like omikuji (a form of drawing lots). The earliest origins of Shinto are lost to history, but it seems to have been established by the late Jomon period. A number of theories about the ancestors of today's Japanese people exist. In addition to the popularly accepted theories of migration from central Asia or Indonesia, some less accepted theories suggest links between modern Shinto practices and ancient Jewish traditions . Most likely, after the arrival of the earliest ancestors of today's Japanese, each tribe and area had its own collection of gods and rituals with no formal relationship between each of the areas. Following the ascendency of the ancestors of today's Imperial family to a position of power among the other groups, their ancestral deities were given prominence over the deities of other groups, though different systems continued to coexist. The introductions of writing in the 5th century and Buddhism in the 6th century had a profound impact on the development of a unified system of Shinto beliefs. In a brief period of time, the Kojiki (The Record of Ancient Things, 712) and the Nihonshoki (The Chronicles of Japan, 720) were written by compiling existing myths and legends into a unified account (see: Japanese mythology). These accounts were written with two purposes in mind. First, the sophistication of the narratives and the introduction of Taoist, Confucian, and Buddhist themes into the narratives were meant to impress the Chinese with the sophistication of the Japanese. The Japanese felt intimidated by the clearly advanced culture of the Chinese and so hoped to produce a work rivaling it. Second, the narratives were meant to shore up support for the legitimacy of the Imperial house, based on its lineage from the Sun Goddess Amaterasu. Much of the area of modern Japan was under only fragmentary control by the Imperial family, and rival ethnic groups (including, perhaps, the ancestors of the Ainu) continued to war against the encroachment of the Japanese. The mythological anthologies, along with other poetry anthologies like the Manyoshu and others, were all meant to impress others with the worthiness of the Imperial family and their divine mandate to rule. With the introduction of Buddhism and its rapid adoption by the court, it was necessary to explain the apparent differences between native Japanese beliefs and Buddhist teachings. One explanation saw the Japanese kami as supernatural beings still caught in the cycle of birth and rebirth. The kami are born, live, die, and are reborn like all other beings in the karmic cycle. However, the kami played a special role in protecting Buddhism and allowing its compassionate teachings to flourish. This explanation was later challenged by Kukai, who saw the kami as different embodiments of the Buddhas themselves. For example, he famously linked Amaterasu, Sun Goddess and ancestor of the Imperial family, with Dainichi Nyorai, a central manifestation of the Buddha, whose name is literally "Great Sun Buddha". In his view, the kami were just Buddhas by another name. Kukai's syncretic view held wide sway up until the end of the Edo period. At that time, there was a renewed interest in "Japanese studies," perhaps as a result of the closed country policy. In the 18th century, various Japanese scholars, in particular Motoori Norinaga (1730–1801), tried to tease apart the "real" Shinto from various foreign influences. The attempt was largely unsuccessful, since as early as the Nihonshoki, parts of the mythology were explicitly borrowed from Chinese doctrines. (For example, the co-creator deities Izanami and Izanagi are explicitly compared to yin and yang.) However, the attempt did set the stage for the arrival of state Shinto, following the Meiji Restoration. Following the Meiji Restoration, Shinto was made the official religion of Japan, and its combination with Buddhism was outlawed. During this period, it was felt that Shinto was needed in order to unify the country around the Emperor as the process of modernization was undertaken with all possible speed. The arrival of large Western gunships and the collapse of the shogunate convinced many that the nation needed to band together if it was going to resist being conquered by outside forces. As a result, Shinto was used as a tool for promoting Emperor (and Empire) worship, and Shinto was exported into conquered territories like Hokkaido and Korea. The era of state Shinto came to an abrupt close with the end of World War II. It appeared that the kami had failed to provide a Divine Wind (kamikaze) to turn back the foreign invaders. Soon after the war, the Emperor even issued a statement renouncing his claims to the status of "living god." In the aftermath of the war, most Japanese came to believe that the hubris of Empire had led to their downfall. Lust for foreign territory blinded their leaders to the importance of their homeland. In the post-war period, numerous "new religions" cropped up, many of them ostensibly based on Shinto, but on the whole Japanese religiosity decreased. Following the war, Shinto has for the most part persisted without the focus on mythology or the divine mandate of the Imperial family. Instead, shrines tend to focus on helping ordinary people gain better fortunes for themselves through maintaining good relations with their ancestors and other kami. Shinto ways of thinking continue to be an important part of the Japanese mindset, though the number of people who identify themselves as religious has suffered a sharp decline. Definition of Shinto Shinto is a difficult religion to classify. On the one hand, it can be seen as merely a highly sophisticated form of animism and may be regarded as a primal religion . On the other hand, Shinto beliefs and ways of thinking are deeply embedded in the subconscious fabric of modern Japanese society. The afterlife is not a primary concern in Shinto, and much more emphasis is placed on fitting into this world, instead of preparing for the next. Shinto has no binding set of dogma, no holiest place for worshippers, no person or kami deemed holiest, and no defined set of prayers. Instead, Shinto is a collection of rituals and methods meant to mediate the relations of living humans to kami. These practices have originated organically in Japan over a span of many centuries and have been influenced by Japan's contact with the religions of other nations, especially China. Notice, for example, that the word Shinto is itself of Chinese origin and that much of the codification of Shinto mythology was done with the explicit aim of answering Chinese cultural influence. Conversely, Shinto had and continues to have an impact on the practice of other religions within Japan. In particular, one could even make a case for discussing it under the heading of Japanese Buddhism, since these two religions have exercised a profound influence on each other throughout Japanese history. Further, the Japanese "New religions" that have emerged since the end of the Second World War have also shown a clear Shinto influence. Some feel Shinto was used as a legitimising ideology during the militaristic phase of Japanese history following the Meiji Restoration. Because Shinto has no absolute source of authority, some feel what was a natural expression of the beliefs of the people was hijacked by radical Nationalists, who desired to unify the Japanese people against the "inferior" people in other nations. Others wonder if the emphasis Shinto places on Japanese exceptionalism made such developments inevitable. Even today, some far right factions within Japanese society want to see a greater emphasis placed on Shinto and increased reverence shown to the Emperor as part of a project to restore Japan to its "rightful place" as the leading nation of the world. However, for most Japanese, Shinto is not about expressing disdain for other nations but expressing one's own love of the natural landscape of Japan and the people and spirits that reside within it. Types of Shinto In order to distinguish between these different focuses of emphasis within Shinto, many feel it is important to separate Shinto into four related types of Shinto expression. - State Shinto was the result of the Meiji dynasty's restoration and the downfall of the shogunate. The Meiji attempted to purify Shinto by abolishing many Buddhist and Confucian ideals; also, the emperor was once again considered divine. After Japan's defeat in World War II, State Shinto was abolished and the emperor was forced to renounce his divine right. - Shrine Shinto is the oldest and most prevalent of the Shinto types. It has always been a part of Japan's history and constitutes the main current of Shinto tradition. - Sect Shinto is comprised of thirteen groups formed during the 19th century. They do not have shrines, but conduct religious activities in meeting halls. Shinto sects include the mountain-worship sects, who focus on worshipping mountains like Mt. Fuji, faith-healing sects, purification sects, Confucian sects, and Revival Shinto sects. - Folk Shinto includes the numerous but fragmented folk beliefs in deities and spirits. Practices include divination, spirit possession, and shamanic healing. Some of their practices come from Taoism, Buddhism, or Confucianism, but some come from ancient local traditions. Characteristics of Shinto The most immediately striking theme in the Shinto religion is a great love and reverence for nature. Thus, a waterfall, the moon, or just an oddly shaped rock might come to be regarded as a kami; so might charismatic persons or more abstract entities like growth and fertility. As time went by, the original nature-worshipping roots of the religion, while never lost entirely, became attenuated and the kami took on more reified and anthropomorphic forms, with a formidable corpus of myth attached to them. (See also: Japanese mythology.) The kami, though, are not transcendent deities in the usual Western and Indian sense of the word—although divine, they are close to us; they inhabit the same world as we do, make the same mistakes as we do, and feel and think the same way as we do. Those who died would automatically be added to the rank of kami regardless of their human doings. (Though it is thought that one can become a ghost under certain circumstances involving unsettled disputes in life.) Belief is not a central aspect in Shinto, and proper observation of ritual is more important than whether one "truly believes" in the ritual. Thus, even those believing other religions may be venerated as kami after death, if there are Shinto believers who wish them to be. Practice and teaching of Shinto Unlike many religions, one does not need to publicly profess belief in Shinto to be a Shintoist. Whenever a child is born in Japan, a local Shinto shrine adds the child's name to a list kept at the shrine and declares him or her "Ujiko", lit. name child. After death an "Ujiko" becomes an "Ujigami", lit. name kami. One may choose to have one's name added to another list when moving and then be listed at both places. Names can be added to the list without consent and regardless of the beliefs of the person added to the list. However, this is not considered an imposition of belief, but a sign of the welcome of the area kami, with the promise of addition to the pantheon of kami after death. Those children who die before addition to the list are called "Mizuko", lit. water child, and believed to cause troubles and plagues. "Mizuko" are often worshipped in a Shinto shrine dedicated to stilling their anger and sadness. These shrines have become more popular with the growth of abortion in modern Japan. Because Shinto has co-existed with Buddhism for well over a millennium, it is very difficult to disentangle Shinto and Buddhist beliefs about the world. One might say that where Buddhism emphasizes the afterlife and ending the cycle of rebirths, Shinto emphasizes this life and finding happiness within it. Though Buddhism and Shinto have very different perspectives on the world, most Japanese do not see any need to reconcile these two very different religions and practice both. Thus, it is common for people to practice Shinto in life yet have a Buddhist funeral. Their different perspectives on the afterlife are seen as complementing each other, and frequently the ritual practice of one will have an origin in the other. Though Shinto has no absolute commandments for its adherents outside of living "a simple and harmonious life with nature and people", there are said to be "Four Affirmations" of the Shinto spirit: - Tradition and the family: The family is seen as the main mechanism by which traditions are preserved. Their main celebrations relate to birth and marriage. - Love of nature: Nature is sacred; to be in contact with nature is to be close to the kami. Natural objects are worshipped as containing sacred spirits. - Physical cleanliness: Followers of Shinto take baths, wash their hands, and rinse out their mouth often. - "Matsuri": Festivals in which the worship and honor is given to the kami. Sin and impurity Shinto does not teach that anything is a sin per se. Rather, certain deeds create a kind of ritual impurity that one should want cleansed merely for one's own peace of mind and good fortune, and not because impurity is wrong in and of itself. Evil and wrong deeds are called kegare (literally, "dirtiness"), and the opposite notion is kiyome (purity). Normal days are called ke (day), and festive days are called hare (sunny, or simply good). Killing living beings should be done with a gratitude and with a reverence for taking life to continue one's own, and it should be kept to a minimum. Modern Japanese continue to place great emphasis on the importance of aisatsu, or ritual phrases and greetings. Before eating, many (though not all) Japanese say "itadakimasu" ("I will humbly receive [this food]") in order to show proper thankfulness to the preparer of the meal in particular and more generally to all those living things that lost their lives to make the meal. Failure to show proper respect can be seen as a sign of pride and lack of concern for others. Such an attitude is looked down upon because it is believed to create problems for all. Those who fail to take into account the feelings of other people and kami will only attract ruin for themselves. The worst expression of such an attitude is the taking of another's life for personal advancement or enjoyment. Those killed without being shown gratitude for their sacrifice will hold urami, (a grudge) and become aragami, a powerful and evil kami that seeks revenge. This same emphasis on the need for cooperation and collaboration can be seen throughout Japanese culture even today. Thus, at modern Japanese companies no action is taken before consensus is reached (even if only superficially) among all parties to a decision. Purification rites are a vital part of Shinto. These may serve to placate any restive kami, for instance when their shrine had to be relocated. Such ceremonies have also been adapted to modern life. For example, a ceremony was held in 1969 to hallow the Apollo 11 mission to the moon, new buildings made in Japan are frequently blessed by a Shinto priest during the groundbreaking ceremony, and many cars made in Japan have been blessed as part of the assembly process. A more personal purification rite is the purification by water. This may involve standing beneath a waterfall or performing ritual ablutions in a river-mouth or in the sea. A third form of purification is avoidance, that is, the taboo placed on certain persons or acts. For example, women were not allowed to climb Mount Fuji until 1868, in the era of the Meiji Restoration. Although this aspect has decreased in recent years, religious Japanese will not use an inauspicious word like "cut" at a wedding, nor will they attend a wedding if they have recently been bereaved. The principal worship of kami is done at public shrines, although home worship at small private shrines (sometimes only a high shelf with a few ritual objects) is also common. It is also possible to worship objects or people while they exist. While a few of the public shrines are elaborate structures, most are small buildings in the characteristic Japanese architectural style. Shrines are commonly fronted by a distinctive Japanese gate (torii) made of two uprights and two crossbars. These gates are there as a part of the barrier to separate our living world and the world the kami live in. There are often two guardian animals placed at each side of the gate and they serve to protect the entrance. There are well over 100,000 of these shrines in operation today, each with its retinue of Shinto priests. Shinto priests often wear a ceremonial robe called a jo-e. Kami are invoked at such important ceremonies as weddings and entry into university. The kami are commonly petitioned for quite earthly benefits; a child, a promotion, a happier life. While one may wish for ill bidding on others, this is believed to be possible only if the target has committed wrongs first, or if one is willing to offer one's life. Though Shinto is popular for these occasions, when it comes to funerals, most Japanese turn to Buddhist ceremonies, since the emphasis in Shinto is on this life and not the next. Almost all festivals in Japan are hosted by local Shinto shrines and these festivals are open to all those that wish to attend. While these could be said to be religious events, Japanese do not regard these events as religious since everyone can attend, regardless of personal beliefs. Shinto's kami are collectively called Yaoyorozu no Kami(八百万の神), lit. eight million kami. The arcane name of eight million, Yaoyorozu is not the exact number, but the expression of infinite number from the time when the concept of infinity did not exist. While such usage has largely disappeared from the common use, until recently there were small shops often in suburb that offered everything from perishable items like foods to magazines and newspapers, even occasionally a bicycle or a car, that was called Yorozu-ya(万屋), lit. 10,000 shop, as to show wide variety of items it offered. The most widely worshipped of all kami is the sun-goddess Amaterasu. However, Japanese do not specifically worship her or call her name to ask for help. Her main shrine is at Ise, but many lesser shrines are dedicated to her. Within the shrine, she is often symbolised by a mirror. Alternatively, the inner sanctum may be empty. This emptiness does not mean non-existence; rather, everything that one sees through the mirror is the embodiment of Amaterasu and every other kami. Until the end of World War II, the Tenno (Emperor) was believed to have been descended from Amaterasu and father of all Japanese, and was therefore a kami on earth (an ikigami or "living kami"); this divine status was popularized during the Meiji restoration. This did not prevent military governors (Shogun) from usurping power, but the emperor was always seen as the true ruler of Japan, even when his rule was only nominal. Although emperor Hirohito renounced his divine status in 1946 under American pressure (Ningen-sengen), the imperial family remains deeply involved in the Shinto ritual that unifies the Japanese nation symbolically. Because Shinto doesn't require a declaration or an enforcement to be worshipped, which is actually "unharmonious" and is something to be avoided, this declaration, while serving political reasons, is religiously meaningless and merely means that the state enforcement has ended. In medieval times, wealthy people would donate horses to shrines, especially when making a request of the god of the shrine (for example, when praying for victory in battle). For smaller favors, giving a picture of a horse became customary, and these ema are popular today. The visitor to a shrine purchases a wooden tablet with a likeness of a horse, or nowadays, something else (a snake, an arrow, even a portrait of Thomas Edison), writes a wish on the tablet, and hangs it at the shrine. In some cases, if the wish comes true, the person hangs another ema at the shrine in gratitude. Cultural effect of Shinto The influence of Shinto on Japanese culture can hardly be overestimated. Although it is now near-impossible to disentangle its influence from that of Buddhism, it is clear that the spirit of being one with nature that gave rise to this religion underlies such typically Japanese arts as flower-arranging (ikebana) and traditional Japanese architecture and garden design. A more explicit link to Shinto is seen in sumo wrestling: the purification of the wrestling arena by the sprinkling of salt and the many other ceremonies that must be performed before a bout can begin are definitely Shinto in origin. It is still very common for Japanese to say, "Itadakimasu" (I humbly partake) before eating, and the Japanese emphasis on proper greetings can be seen as a continuation of the ancient Shinto belief in kotodama (words with a magical effect on the world). Many Japanese cultural customs, like using wooden chopsticks and removing shoes before entering a building, have their origin in Shinto beliefs and practices. Also, a number of other Japanese religions, including Tenrikyo, have originated from or been influenced by Shinto. - Atsuta Shrine, Nagoya, Aichi, shrine to the Imperial sword Kusanagi - Heian Jingu (Kyoto), dedicated to Emperors Kammu and Kōmei - Ise Shrine (Ise), dedicated to Amaterasu - Itsukushima Shrine, Hiroshima prefecture - Iwashimizu Shrine, Yawata, Kyoto - Izumo Shrine (Izumo) - Kasuga Shrine, Nara - Katori Shrine , Chiba Prefecture - Kumano Shrines , Wakayama Prefecture - Meiji Shrine (Tokyo), the shrine of Emperor Meiji - Nikko Toshogu, Nikko, Tochigi Prefecture - Tsurugaoka Hachiman Shrine, Kamakura, Kanagawa - Usa Hachiman Shrine , Oita Prefecture - Yasukuni Shrine (Tokyo), controversial shrine dedicated to the 'peace of the nation' and seen by some as a symbol of Japan's militaristic past See also List of Shinto shrines - Culture of Japan - History of Japan - Japanese Buddhism - Japanese mythology - Japanese nationalism - Jinja (shrine) - List of Shinto shrines - Religions of Japan - Shinto music - Shinto - a Philosophical Introduction by Timothy Takemoto - Shinto struggles to find a place in postwar society - a basic introduction to Shinto by Eric Talmadge for the Japan Times The contents of this article is licensed from www.wikipedia.org under the GNU Free Documentation License. Click here to see the transparent copy and copyright details

http://all-science-fair-projects.com/science_fair_projects_encyclopedia/Shinto

Spain's New World Exploration Time Line Students Create a Time Line of Spain's Exploration of the Americas Students use the time line template to outline Spain's exploration of the Americas. Students identify important dates in the history of Spanish exploration by major explorers such as Columbus, Balboa, Ponce de Leon, Magellan, Cortes, Pizarro, Desoto and Coronado. Students add the years of exploration and the explorer's name in the text boxes. Students add brief information about the explorations in the textboxes. Students connect the text boxes to the appropriate dates on the time line with arrows. This activity integrates technology into the History - Social Sciences curriculum ( Age of Exploration ). Grade 4 and up 1) Research the internet and textbooks to identify important dates of Spain's exploration of the Americas by major explorers. 2) Enter the dates of exploration, the explorer's name, and a brief fact about the exploration into the textboxes. 3) Use arrows or lines to connect the text boxes to the dates on the time line. Spain's New World Exploration Time Line Template:

http://oakdome.com/k5/lesson-plans/word/new-world-exploration-time-line-spain.php

and aerospace engineers involves the theory of beams. A fundamental concept in beam theory is the idea of a centroid, or more generally, the center of gravity. Given a strangely-shaped region cut out of stiff cardboard, students can readily understand the idea of center of mass: it's the point where the region would exactly balance on the tip of a pencil. A more difficult problem is how to compute the location of that point, especially if the material is not symmetric and is not of uniform density. As soon as the multivariable calculus students are introduced to double integrals, they are immediately asked to think about the geometric ideas associated with numerically approximating the area of a planar region. They use Maple to help them explore the idea of superimposing a grid on a region and adding up the area of all squares that are "inside" the region. The region we give them is a model of the state of Minnesota (see Figure 1); the location of the state boundary was derived from a cartographic database. We call this the "state of Minnesota lab". Once the students have learned and extended some of these ideas, they are asked to compute the centroid of the state, and then to consider what happens to the "center" of Minnesota if we do not assume a uniform density distribution. Specifically, we create a model that takes into account a real geological feature: the Great Iron Range of northern Minnesota. The model represents the Iron Range by assuming a density distribution which is highest in a roughly east-west swath across the upper portion of the state. The students discover that if we assume such a density distribution, then the state's "geological" center of mass is further north-east than if we assume a uniform distribution. The students then look at a model of a population distribution in which the bulk of the state's population lies in the eastern portion of the state, and in which the state's population decreases from south to north. They conclude that the demographic center of this model is just northwest of the Twin Cities. The students are now equipped to take a density function for any distribution of objects (people, soybeans, income, voters, etc.) and find the center of mass for that distribution, but they also learn more from this lab than just how to find a center of mass. They learn to question the accuracy of a model and to interpret a model's predictions. They learn that some problems have no closed-form analytic solutions. They learn both theoretical and numerical techniques of integration, and why each is important. Finally, they learn these concepts better with by using technological tools, because of the rich exploration that technology permits. Author: Frederick J. Wicklin <[email protected]> Comments to: [email protected] Created: Fri Aug 18 1995 --- Last modified: Jul 21 1996 Copyright © 1995-1996 by The Geometry Center All rights reserved.

http://geom.math.uiuc.edu/docs/outreach/PWCEE.95/MN.html

communities and governments that work for and support human rights for children provide ethnic and national knowledge and roots for their children. They name their children, and help them acquire a sense of belonging in their family, nation, and world. Through this belonging, their children become invested in the positive development of their family and nation (CRC Articles 7, 8). Rights: Children have the right to a name; to a nationality. Responsibilities: Children are responsible for respecting the rights of those who live in or come from other countries; standing up for their own and other's rights to a name, nationality, and other indicators of identity; working toward the positive development of their * gain a beginning understanding of the * increase their understanding about their nationality, race, ethnicity, gender, and life role; * increase respect for their own and others' * share information with their children that helps them understand their heritage; * teach and role model, according to their child's evolving capacity to learn, responsibilities pertaining to their name and their nationality; * explore ways to augment children's self-concept with knowledge about their name and heritage; * explore the formation of our collective - Flags of the World chart; - Native Cultures flag chart: - Colored paper, scissors, glue (red, white and blue), and markers; - Rice, scoops, cups and spoons (You might provide a variety of rice, so participants can see the differences); - Red and blue paint, white paper, paint brushes, star stamp; - Construction paper American flags with instructions on them (see Parent/Child Interactive Activities, Name and - Chart paper and markers; - Name cards (from early childhood - Wee Sing Around The World audiotape; - Raffi's One Light, One Sun audiotape, ;Like Me and You; song; - Colorful, thin-point markers; - Extra copies of the Convention on the Rights of the Child. Greet as usual. Make sure everyone gets a name tag. Parent/Child Interactive Activities - 1. FLAGS (CREATIVE EXPRESSION) The flag from our country symbolizes the nation that we call our homeland. - Families make flags of the place (country, tribe, area, region) from where their ancestors came. - Supply charts which show various flags of nations and tribes. Rice is a food with which a majority of the world's people are familiar. - Place uncooked rice of several varieties and scoops, etc. into the sensory table. Suggest parents help children in sorting and naming the varieties. - 3. AMERICAN FLAG (COOPERATIVE ART) The American flag is the symbol for the United States of America (USA). The 50 stars represent the 50 states. The 13 stripes represent the original 13 colonies. - Provide a star stamp and red and white paint. Parent and children will make the American flag together by making red stripes and stamping stars onto blue paper in top left hand corner. - 4. NAME AND NATION WALK (SMALL AND LARGE MUSCLES) Provides a vehicle for parentchild discussion about name and nationality. This discussion is preparatory for the parent discussion topic of the day. - Use flag shapes and write instructions on them. Put the flags around the room. Have parents and children walk around the room, read them, and do the actions requested on the flag. Name and Nation Walk Preparation: Make flags with some, or all, of the below instructions on them. - Tell your child your full name, ask him or her to say their full name. - Tell your child whom he or she is - Tell your child the meaning of his or her name. - Tell what you know about the ethnic origin of your child's name. - Finger spell your name to your child (Use the American Sign Language Finger Spelling Chart). Finger spell your child's name, and encourage children to finger spell their own names. - Tell your child what your nickname is and how you got it. Tell your child how he or she got their nickname or why they don't have one. - Tell your child what country your ancestors came from. - Show your child the flags of all the countries your ancestors came from. - 5. COLOR NAME (CREATIVE EXPRESSION) Gives a visual way to celebrate names and the value of each person as an individual. It also reminds children that they are part of this country. - This activity uses the cue card from Session 3. Fold an 8 ; inch piece of paper in half. Then - write your child's name above the crease with different colors of glue. Red, white and blue glue are provided. Fold paper again and pat down, open paper and have your child sprinkle glitter to create a mirror image of their name. Using the red, white and blue glue will remind children of the flag of the United - 6. BOOK CORNER (LANGUAGE) - Everybody Cooks Rice, by Ann - A Flag for Our Country, by - Families Are Different, by - I Hate English, by Ellen Levine - Everybody Cooks Rice, Norah - 1. Transition: Early childhood teacher speaks to each child, and/or touches them on the shoulder and reminds them that circle time will begin soon. After connecting with each child, the teacher begins a gathering song. - 2. I'm happy to see all of you!; Sing a get acquainted song of your choice, or sing, ;Shake Hands With Friends and Say Hello; and the ;Name - 3. Today, our theme is name and nationality. Let's sing a song to recognize all the children here today.; - 4. Explain: ;Your name is special even if you know someone with a name like yours or the same as your name, your parents gave you a name that they thought was just right for who you are. Your name is as precious as a jewel. So is every other person's name. Names are precious and need to be protected. It's very important that no one ever makes fun of someone's name. Later we will talk more about names. - ;Parents, as we go around this circle will you shout out your child's complete name? You say the first, middle and last names, and then we will sing this song using their first name only.; - 5. Sing: If your name is ______ stand up tall.; The teacher introduces the ;Name Game.; Invite everyone to stand up. Invite parents to help their children point to the person being named. It goes like this: a. The teacher selects a child to begin by singing: Ann, Ann, look at everyone Point to Sue and then you're done. b. After the child points to Sue, she/he sits down, and the teacher continues: Sue, Sue, look at everyone Point to Bill and then you're done.; c. Continue in this way until all the children are named. If your group is small enough (eight or less), name parents as - 6. Sing ;Shake Hands With Friends; again, and ask that participants say, ;Hello, _____ (child's name), I hope that we can be friends today,; as they are singing the song. In other words, participants use people's names with their handshake. - 7. I'd like to go around the circle one more time and have each child, with help from their parent, tell us which countries you or your ancestors came from. Here in the United States there are people from all over the world. Let's find out which countries are represented in our class.; - Begin by stating which country/ies you or your parents, grandparents or great grandparents came from. Then the child to the right or left tells about his or her ancestry, and so on, around the circle. After the last child/parent has shared, thank everyone for sharing - 8. Sing: ;The More We Get Together,; (using sign language signs, if possible) and ;This Land is Your Land.; - 9. Close with: ;This Little Light of Mine.; NOTE: Adults may have to help children begin this game. As children get comfortable, they will not be shy. Separate learning time Children's Learning Circle Session - 1. Invite children to the circle with a gathering song. - 2. Teacher says: ;Remember when we talked about names in the big circle today? Let's remember everybody's name again. Go around the circle and as a group, say everyone's name together. - 3. Sing: ;The Name Chant; or ;Everybody Stand up Tall.; - 4. Ask the children if they can remember what country their ancestors came from. Ask the children what country they live in now. - 5. Share the American flag with the children. Count the stripes and stars together. - 6. ;The American Flag is the symbol for our country, the United States of America. Sometimes it is called America, or the USA. Those are different names for the same country. There are fifty stars on our flag. Each star represents, or stands for a state in our country. The state we live in is _________. There are thirteen stripes on our flag. Each stripe counts for one of the thirteen colonies that were the original states when our country was born.; - 7. Sing: ;This Land is Your Land; or ;This Little Light of Mine.; - 8. ;Now, we have a color flag game to help us learn about the colors in our flag.; (From Hap Palmer record: Learning Basic Skills Through Music AR 514 Vol 1. Original words and music by Hap Palmer.) Hand out red, blue, green and yellow flags to all the children. - ;Let's listen to what this song says and follow the directions. It will tell us to stand up or sit down. Let's all try that now. We will need to listen very carefully. Look at what color your flag is. When you hear your color name, then stand up or sit down according to what the song says.; - 9. Ask for favorite songs from the children and sing them. - 10. Close the circle with: ;The More We Get Together.; Parent Education Session 4 Preparation: ;Name and nationality.; Write this topic title on chart paper or chalk board. As parents enter the room, have the Wee Sing Around the World audiotape playing. Write on newsprint, ;SIGN IN, PLEASE! Please write your entire name on this newsprint.; Provide thin-pointed markers for participants to write their names. As soon as everyone is assembled, turn off the tape. - 1. Greeting: ;Shalom! Bonjour! Buenas Tardes! G'day! Guten Tag! Cio! Nyob zoo! We are ready to begin. Today our topic is 'Being and Belonging.'; - 2. Names Group instructions: a. Say your entire name as it is written on the newsprint, as well as the full name of your child. b. If you know how to say hello in one of the languages of your family's origin, c. State something you believe about names. - 3. Discussion and questions: How/why does your name hold importance to you? What do our names give us? What does our language of origin give us? - 4. Name art: Some of you did Name Art cards with your child today. Please share your creation and your ;name story; with the group, if you have one. For example, you might tell us the significance or meaning of your name, whom you are named after, and so on. You might also tell us about your choice of colors (if you provided - 5. Think about the activities you and your child just worked on together. What did you notice about your child's interest in or reaction to one of the activities? (flags, color name art, Parents focus on children's feelings or discussion during interaction time activities. Parents interpret children's - 6. How does a name relate to Parents make connections between who they are known as and how they know themselves. Names identify who we are to others. - 7. Nationality: ;We are addressing the 'Right to a name and nationality,' during this session. We just talked about our own names and how our names may affect us. Now let's talk about how our nationality impacts our lives. We often take our nationality for granted, rather than recognize how powerful an impact it has on how we see ourselves, how we see each other, and how we see the world. For example, the Pledge of Allegiance is a defining document for us in the United States. Does everyone know it?; If not, recite it for them. - ;Does anyone want to share their thoughts about this pledge? . . . How does this pledge describe us?. . . What does it say about our nationality? . . . How do we feel about that?; Invite open discussion. Remind participants, if necessary, to appreciate each person's - 8. Brainstorming activity: Collective a. Write the word ;nationality; on the chalkboard and give parents time to reflect on its meaning. Chart their b. Together identify as many things as possible that we share because we live in the United States. Make a list. When finished say, ;This list tells us about our collective identity. Are these things that make you proud that you live in this c. List things people wish were not part of our collective identity. - 9. Brainstorming activity: Standing up for one's country Ways we typically stand up for the country we consider our homeland. (List.) Ways we listed that we can use while also honoring the rights of people in Put a star by the few that meet this criteria. For example, we may stand up for our country by going to war to protect her. However, this criteria would not receive a star because it is not good for other countries. When we stand up for our country through peaceful means we can show our support for our country without showing disrespect for other countries. "How we can stand up for our country in ways that teach our children about compassion, embracing differences, peacemaking, and so on, and generally role model what we want them to learn?" "How can we impact our collective identity and make a statement about who we want to be in this country, while standing up for our country?" - 10. Summary: ;Human beings have a basic need to belong. They must know themselves and how they fit into the world. They must know who they are and to whom and what they belong, or of what they are a part. For these reasons, having a name and a nationality are basic human rights. When these rights are honored, children can know themselves and their country. Through developing a deeper understanding of their name and their nationality, they can go beyond blind acceptance of that identity and learn to question it. This questioning is part of our identity in this country. - ;Our children, and we, their parents, can make a difference for our homeland by standing up for what is right, knowing that part of our collective identity is honoring liberty and justice for all. In this way we increase our respect for ourselves, and we impact the collective identity in positive ways.; a. What are those things you hope your child values about his or her family or about his or her country? b. While listening to Raffi's "Like Me and You" song: Reflect on the music and words to this Record some thoughts about your child's name and your family's national ties.

http://www1.umn.edu/humanrts/edumat/hreduseries/rrr/sess4.html

In this three-day simulation lesson, students explain the steps taken from party formation to national election. Harnessing skills gained from the Electoral Process lesson, students will act out the campaigning and voting process by simulating a real election in their own classroom. Explore the evolution of voting rights in the Unites States through an interactive PowerPoint presentation highlighting landmark changes. Following the presentation and class discussion, students apply the new knowledge of voting legislation to individual scenarios through a class activity. voting; 15th Amendment; 19th Amendment; 23rd Amendment; 26th Amendment; Voting Rights Act; Poll Tax; Poll Tests; Grandfather Clause; suffrage; voting age; Susan B. Anthony; Elizabeth Cady Stanton; literacy test; civil rights In this lesson, students evaluate hypothetical candidates by establishing and applying their own criteria for selecting public officials. Through a variety of activities, students assess candidates based on their qualifications, experience, campaign speeches and campaign materials. Students track campaign promises, explore voting records and evaluate the legitimacy of information resources. The role of the media, fundraising and opinion polls in the electoral process is also discussed. Students will learn about the essential characteristics of state government including the duties and functions associated with the legislative, executive and judicial branches. This lesson will reveal the impact that agencies amd commissions have on students' lives and illustrate the law-making process at the state level. State government; governor; state legislature; state senate; state house of representatives; state courts; passing bills; committees In this lesson, students will compare the executive, legislative and judicial branches at all levels of government. Students will learn why powers and obligations are distributed between the levels of government. Students grasp the nuances of diplomacy through this interactive lesson. They are called to decide which diplomacy tools work best in different situations. Students will develop an understanding of negotiation, sanctions, and other elements used in diplomatic relationships. foreign policy; isolationism; internationalism; national interest; diplomatic strategy; aid; sanctions; military force; mediation; negotiations; treaty Countries often work together to solve problems and fall into conflict when problems cannot be resolved. After learning about motivations and conditions that lead to action (or inaction), students analyze examples of international conflict and cooperation. conflict; cooperation; conditions; motivations; actions; international cooperation; international conflict Economic, cultural, and military influence are all critical in developing spheres of influence. Students explore international authority by following a Cold War case study, which will encourage better understanding of international persuasion. sphere of influence; containment; capitalism; communism; propaganda; Truman Doctrine; Cold War; hot war; NATO; Warsaw Pact; Marshall Plan Students compare the basic structure of several different international organizations before categorizing their work. Students also examine the local and global impact of international organizations. international organizations; nongovernmental organizations; NGO; intergovernmental organization; World Health Organization; Red Cross/Red Crescent; World Bank; NATO; European Union; United Nations; UNICEF Directions for Democracy has been updated and replaced by Anatomy of the Constitution. This lesson introduces the Constitution of the United States. Students will interpret the intentions of the Preamble, explain the organization of the U.S. government, and identify the rights protected in the Bill of Rights. constitution; Preamble; Bill of Rights; articles; amendments; ratify; compromise; Anti-federalist; Federalist; individual freedoms

http://www.icivics.org/teachers/lesson-plans?page=3

COMMON AMATEUR ASTRONOMY TERMS The user is assisted in exploiting the formulas found in this book through usage notes, definitions, and examples provided throughout the Astro Functions and Methods. This sheet lists some common terms and concepts used throughout the work. · UT is Universal Time, which is the standard time at the prime meridian (0-degrees longitude) running through Greenwich England. UT times are given on a 24-hour clock. In the Americas a number of hours must be added to local time to calculate UT. In the continental USA the standard corrections are +5 (Eastern), +6 (Central), +7 (Mountain), and +8 (Pacific) hours. Add one hour less when daylight savings is in effect. Note that, if this addition causes the time to pass midnight (exceeds 24-hours) you must increment your calendar date. For instance, Central Standard Time (CST) is 6-hours behind UT, if it were 8:44 PM CST on May 20th you would determine UT as follows: · The celestial equator is a circle of reference created by an extension of the Earth's equator into space. For an observer standing on the equator, it would run dead east-west through the zenith -- the highest point in the sky. Observers at the poles would have the celestial equator running along the horizon. The ecliptic is another reference circle, created by using the plane of the Earth's orbit. The path taken by the sun across the sky traces a section of the ecliptic each day. The moon and planets move some degrees north or south of this circle. The ecliptic and celestial equator would be the same circle if the Earth's axis of rotation were perpendicular to its orbit. But, the planet is tilted: so these circles intersect each other at the two equinoxes and form an angle called the obliquity of the ecliptic. The vernal equinox is the intersection point that the sun reaches in spring and is used as the starting point for measuring angular distances along the ecliptic or equator. · Right ascension (RA) and declination (DEC) form the celestial equatorial system of measure, that uses the vernal equinox and celestial equator as starting points. It is similar to the system of longitude and latitude on the surface of the globe. RA is measured eastward along the celestial equator from the vernal equinox. It is given in units of hours, which correspond to 15 degrees of arc. In this way, 24 hours of RA equal 360 degrees of arc (24 x 15 = 360). These hours are subdivided into minutes and seconds, just like the hours on a clock (see the Mean Solar Day to Sidereal Day function at the end of Basic Conversions on why you can't use your watch to measure off RA). For purposes of calculation these hour:minute:second of position are first converted into degrees of arc. Declination is measured in degrees from the celestial equator (0-degrees) north (+) and south (-) to the celestial poles, which reside at +/-90 degrees of declination and coincide with the rotation axis of the planet. · The next most common positional system encountered by the amateur astronomer is celestial ecliptic. It shares with the RA and DEC system in the use of the vernal equinox as a starting point for positive eastward measurement. Measurement, though, is along or perpendicular to the circle of the ecliptic. The angular distance north or south of this baseline is ecliptic or celestial latitude. It runs up to ±90-degrees ending at the ecliptic poles. Ecliptic or celestial longitude is the angular distance eastward along the ecliptic from the vernal equinox point. Unlike RA and DEC, both measurements are usually given in degrees. · Gravitational effects on the Earth, mainly from the Sun and Moon, cause the equinox points to shift along the celestial equator. A long term effect, known as precession, causes the celestial poles to rotate around the ecliptic poles in a cycle of 26,000 years. This has the unfortunate effect of allowing celestial coordinates for an object to change over time. Therefore, all such coordinates are given in terms of a date epoch. Currently the standard epoch is known as J2000.0. This is equivalent to noontime on the first day of the year 2000. You may also see epochs for quarter and midyears: 1991.25, 1999.5, etc., as well as instantaneous epochs. Astro Utilities Electronic Book Copyright © 1999 Pietro Carboni. All rights reserved.

http://www.pietro.org/Astro_Util_StaticDemo/CommonAstronomyTerms.htm

The roots of the Inca culture date back to around the year 1200 A.D., but the period of the Inca empire started only in the 15th century. In under a hundred years, the Incas created a state that extended from the border zone of Columbia and Ecuador to the central parts of Chile. Circa nine million people lived in the area. The short golden age of the Inca state ended when the Spanish troops led by Francisco Pizarro arrived in Peru in the year 1532 and subjugated the central areas of the Inca state. The success of the Inca state was largely based on an extensive network of roads and a functioning system of administration. The road network made it possible to move troops quickly, and messengers placed along it could deliver the message from one end of the state to another very fast. Also the ritualistic exchange of gifts between the Inca ruler and local chieftains was important. With the help of gifts and strategic marriages the Inca ruler made even the remotest chieftains his subjects and relatives. With the capital Cusco as its centre, the Inca empire was divided into four administrative areas. The state was ruled by the Inca king. On the next level of the hierarchy were the prominent persons of Cusco and the leaders of the bigger ethnic groups subjugated by the Inca state. Numerically, the population was divided into units of 10 000, 1 000, 100 and 10 of men of working age and their families. The units had their own leaders. The numerical division of the male population was above all needed for the organization of the system of work taxation. The main deity of the Inca people was the sun god Inti, besides which creator god Viracocha and moon, stars, lightning and rainbow were worshipped. Mainly llamas were sacrificed to gods and numerous sacred places of nature in religious rituals. Human sacrifices were rather rare. A man holding a quipu (knotted string). Quipus were used to prepare tax and inventory lists, and perhaps also for the registration of historical events. There was a specialist profession for the knotting and interpreting of the quipus. The drawing is by Guaman Poma, an Indian-born chronicler. An aryballo jar with a long spout and pair of handles is one of the most typical Inca vessel shapes. Aryballos of various sizes are known from small to huge, over 1 m tall. The aryballos were probably used to hold maize, corn beer, and water. © Museo Nacional de Arqueología, Antropología e Historia del Perú - Instituto Nacional de Cultura del Perú, Lima, Peru (Cat. 291) Machu Picchu, located a few days' journey from the Inca capital Cusco, is the most famous of the Inca sites and one of the best preserved. © Antti Korpisaari The preserved stone walls of Coricancha, the main temple of Cusco, are among the most impressive examples of the world-famous Inca architecture. © Antti Korpisaari

http://www.tampere.fi/ekstrat/taidemuseo/arkisto/peru/800/inka_en.htm

Temporal range: Upper Jurassic–Upper Cretaceous |Deinonychus (large) and Buitreraptor (small) at Field Museum of Natural History, Chicago Matthew & Brown, 1922 Matthew & Brown, 1922 Dromaeosaurs had a world-wide distribution. They first appeared in the Middle Jurassic, 167 million years ago (mya). The larger types are not found before about 124 mya in the Lower Cretaceous, and are only found in the northern hemisphere. They survived until the end of the Cretaceous, 65.5 mya at the K/T extinction event. The presence of dromaeosaurs in the Middle Jurassic is shown by isolated fossil teeth, though no dromaeosaurid body fossils have been found this early. The Dromaeosaur skeleton suggests they were active, fast-running, and closely related to birds. Dromaeosaurs, like most other theropods, had an S-curved neck, and their trunk was relatively short and deep. Like other maniraptorans, they had long arms, and relatively large hands with three long fingers ending in large claws. The dromaeosaur hip structure featured a large 'pubic boot' (part of the hip) projecting beneath the base of the tail. These hip bones were places muscle and tendons attached. Dromaeosaur feet bore a large, recurved claw on the second toe. Their tails were slender, and mainly used as a counter-balance. Some, and probably all, dromaeosaurs were covered in feathers, including large, vaned, wing and tail feathers. This development, first proposed in the mid-late 1980s, was confirmed by fossil discoveries in 1999. It is a significant change in the way dromaeosaurs have been depicted in art and movie. Like other theropods, dromaeosaurs were bipedal; that is, they walked on their hind legs. Whereas other theropods walked with three toes on the ground, fossilized footprint tracks show that most dromaeosaurs held the second toe off the ground, with only the third and fourth toes bearing the weight of the animal. The enlarged second toe bore an unusually large, curved sickle-shaped claw. It is thought to have been used in capturing prey and, in the smaller species, climbing trees. One dromaeosaur species, Balaur bondoc, had a first toe which was highly modified in parallel with the second. Both the first and second toes on each foot of B. bondoc were also held retracted and bore enlarged, sickle-shaped claws. Dromaeosaurs had long tails. Most of the tail vertebrae bear bony, rod-like extensions, as well as bony tendons in some species. In his study of Deinonychus, Ostrom proposed that these features stiffened the tail so that it could only flex at the base, and the whole tail would then move as a single, rigid, lever. However, one well-preserved specimen of Velociraptor mongoliensis (IGM 100/986) has an articulated tail skeleton which is curved horizontally in an S-shape. This suggests that, in life, the tail could bend from side to side with some degree of flexibility. It has been proposed that this tail was used as a stabilizer and/or counterweight while running. In Microraptor, an elongate diamond-shaped fan of feathers is preserved on the end of the tail. This might have been used as a stabilizer and rudder during gliding. Dromaeosaurs were small to medium-sized dinosaurs, ranging from about 0.7 meters in length (2.3 ft, in the case of Mahakala) to over 6m (20 ft, in Utahraptor and Achillobator). Some were larger; undescribed specimens of Utahraptor in BYU collections may have been up to 11 m (36 ft) long. Large size appears to have evolved at least twice among dromaeosaurs. A possible third lineage of giant dromaeosaurs is represented by isolated teeth found on the Isle of Wight, England. The teeth belong to a Utahraptor-sized animal but are more similar in shape to the teeth of velociraptorines. Mahakala is both the most primitive dromaeosaur ever described and the smallest. This evidence, and Microraptor and the troodontid Anchiornis, suggests the common ancestor of dromaeosaurs, troodontids, and birds – the 'ancestral paravian' – may have been very small, at around 65 cm in length and 600 to 700 grams of mass. Fossils show that dromaeosaurs were covered in feathers. Some fossils preserve long feathers on the hands and arms (remiges) and tail (rectrices), as well as shorter, down-like feathers covering the body. Other fossils, which do not preserve actual impressions of feathers, still preserve the associated bumps on the forearm bones where long wing feathers would have attached in life. Overall, this feather pattern looks very much like Archaeopteryx. Many other dromaeosaur fossils have been found with feathers covering their bodies, some with fully-developed feathered wings. Microraptor even shows evidence of a second pair of wings on the hind legs. While direct feather impressions are only possible in fine-grained sediments, some fossils found in coarser rocks show evidence of feathers by the presence of quill knobs, the attachment points for wing feathers possessed by some birds. The dromaeosaurids Rahonavis and Velociraptor have both been found with quill knobs, showing that these forms had feathers despite no impressions having been found. In light of this, it is most likely that even the larger ground-dwelling dromaeosaurs bore feathers, since even flightless birds today retain most of their plumage, and relatively large dromaeosaurs, like Velociraptor, are known to have had feathers. Claw function [change] There is some discussion about the function of the enlarged "sickle claw" on the second toe. When John Ostrom described it for Deinonychus in 1969, he interpreted the claw as a blade-like slashing weapon, much like the canines of some saber-toothed cats, used with powerful kicks to cut into prey. Adams (1987) suggested that the talon was used to disembowel large ceratopsian dinosaurs. The interpretation of the sickle claw as a killing weapon applied to all dromaeosaurs. In Manning's interpretation, the second toe claw would be used as a climbing aid when subduing bigger prey and also as stabbing weapon. Ostrom compared Deinonychus to the ostrich and cassowary. He noted that the bird species can inflict serious injury with the large claw on the second toe. The cassowary has claws up to 125 millimetres (4.9 in) long. The seriema also has an enlarged second toe claw, and uses it to tear apart small prey items for swallowing. The Manning team also compared the curvature of the dromaeosarid "sickle claw" on the foot with curvature in modern birds and mammals. Previous studies had shown that the amount of curvature in a claw corresponded to what lifestyle the animal has: animals with strongly curved claws of a certain shape tend to be climbers, while straighter claws indicate ground-dwelling lifestyles. The sickle-claws of the dromaeosaurid Deinonychus have a curvature of 160 degrees, well within the range of climbing animals. The forelimb claws they studied also fell within the climbing range of curvature. Paleontologist Peter Mackovicky said that small, primitive dromaeosaurids (such as Microraptor) were likely to have been tree-climbers, but that climbing did not explain why later, gigantic dromaeosaurids such as Achillobator retained highly curved claws when they were too large to have climbed trees. Group behavior [change] Deinonychus fossils have been uncovered in small groups near the remains of the herbivore Tenontosaurus, a larger ornithischian dinosaur. This had been interpreted as evidence that these dromaeosaurs hunted in coordinated packs like some modern mammals. However, not all paleontologists found the evidence conclusive, and a study in 2007 by Roach and Brinkman suggests that the Deinonychus may have actually displayed a disorganized mobbing behavior. The first known extensive dromaeosaur trackway was found in Shandong, China. The trackway, (made by a large, Achillobator-sized species), showed the sickle-claw was held off the ground. Six individuals of about equal size moved together along a shoreline. The individuals were spaced about one meter apart, and retained the same direction of travel, walking at a fairly slow pace. The trackways are evidence that some species of dromaeosaurs lived in groups. While the trackways clearly do not represent hunting behavior, the idea that groups of dromaeosaurs may have hunted together can not be ruled out. Flying and gliding [change] The ability to fly or glide has been suggested for at least two dromaeosaur genera. The first, Rahonavis ostromi was originally classified as avian bird, but found to be a dromaeosaur in later studies. It may have been capable of powered flight. The forelimbs of Rahonavis were more powerfully built than Archaeopteryx, and show evidence that they bore strong ligament attachments necessary for flapping flight. Luis Chiappe concluded that, given these adaptations, Rahonavis could probably fly but would have been more clumsy in the air than modern birds. Another species of dromaeosaur, Microraptor gui, may have been capable of gliding using its well-developed wings on both the fore and hind limbs. A 2005 study by Sankar Chatterjee suggested that the wings of Microraptor functioned like a split-level "biplane", and that it likely employed a style of gliding, in which it would launch from a perch and swoop downward in a 'U' shaped curve, then lift again to land on another tree, with the tail and hind wings helping to control its position and speed. Chatterjee also found that Microraptor had the basic requirements to sustain level powered flight in addition to gliding. Relationship with birds [change] - Further information: Origin of birds Mark Norell and colleagues analyzed a survey of coelurosaur fossils and suggested that dromaeosaurs were most closely related to birds, with troodontids as a more distant outgroup. In 2002, Hwang and colleagues suggested that birds (avialans) were better thought of as cousins to the dromaeosaurids and troodontids. The current consensus among paleontologists agrees with Hwang that dromaeosaurids are most closely related to the troodontids, and together with the troodontids form the clade Deinonychosauria. Deinonychosaurians in turn are the sister taxon to avialans, and therefore the closest relatives of avialan birds. A consensus of paleontologists has concluded that there is not yet enough evidence to say whether any dromaeosaurs could fly or glide, or whether they evolved from ancestors that could. Other ideas [change] In 2002, Hwang et al. found that Microraptor was the most primitive dromaeosaur. Xu and colleagues in 2003 cited the basal position of Microraptor, along with feather and wing features, as evidence that the ancestral dromaeosaur could glide. In that case the larger dromaeosaurs would be secondarily terrestrial—having lost the ability to glide later in their evolutionary history. A few researchers, like Larry Martin, believe that dromaeosaurs, along with all maniraptorans are not dinosaurs at all. Martin asserted for decades that birds were unrelated to maniraptorans, but in 2004 he changed his position, and now he agrees that the two are the closest of relatives. Martin believes that maniraptorans are secondarily flightless birds, and that birds evolved from non–dinosaurian archosaurs, so that most of the species formerly called theropods would now not even be classified as dinosaurs. A challenge to all of these alternative scenarios came when Turner and colleagues in 2007 described a new dromaeosaur, Mahakala, which they found to be the most basal and most primitive member of the Dromaeosauridae, more primitive than Microraptor. Mahakala had short arms and no ability to glide. Turner et al. also inferred that flight evolved only in the Avialae, and these two points suggested that the ancestral dromaeosaurid could not glide or fly. Based on this cladistic analysis, Mahakala suggests that the ancestral condition for dromaeosaurids is non-flying. Phylogeny and taxonomy [change] Dromaeosauridae was first defined as a clade by Paul Sereno in 1998, as the most inclusive natural group containing Dromaeosaurus but not Troodon, Ornithomimus or Passer. The various "subfamilies" have also been re-defined as clades. The subfamilies of Dromaeosauridae is not yet settled. Mahakala, for example, the most primitive dromaeosaur in terms of structure, falls outside any named sub-group. The most basal subfamily of dromaeosaurs is often found to be the Unenlagiinae. All known dromaeosaur skin impressions come from this group, and all show an extensive covering of feathers and well-developed wings. Some species may have been capable of active flight. The following classification of the various genera of dromaeosaurids is based on studies by Sereno (2005), Senter (2004), Makovicky et al. (2005), Norell et al. (2006), and Turner et al. (2007). - Family Dromaeosauridae - Subfamily Dromaeosaurinae - Subfamily Microraptorinae - Subfamily Saurornitholestinae - Subfamily Unenlagiinae - Subfamily Velociraptorinae Velociraptorinae is a subfamily of the Dromaeosauridae. It existed from the Late Jurassic period to the end of the Cretaceous period. The Velociraptorinae consists of Velociraptor, Deinonychus, Tsaagan, Saurornitholestes, and Balaur. Teeth belonging to a giant velociraptorine the size of Utahraptor have also been reported from the Isle of Wight, England. In general, velociraptorines tend to have longer, narrower jaws and more slender bodily proportions than the dromaeosaurines. In popular culture [change] A Deinonychus appears in John Brosnan's 1984 novel Carnosaur and its movie adaptation though the book itself received little press attention. Velociraptor, a dromaeosaur, got attention after it was featured in the 1993 Steven Spielberg movie Jurassic Park. 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Peabody Museum of Natural History Bulletin 30: 1–165. - Norell, Mark A.; & Makovicky, Peter J. (1999). "Important features of the dromaeosaurid skeleton II: information from newly collected specimens of Velociraptor mongoliensis". American Museum Novitates 3282: 1–45. http://hdl.handle.net/2246/3025. - Chatterjee S. and Templin R.J. 2007. Biplane wing planform and flight performance of the feathered dinosaur Microraptor gui. Proceedings of the National Academy of Sciences, 104(5): 1576–1580. - Hwang S.H. Norell M.A. Ji Q. and Gao K. 2002. "New Specimens of Microraptor zhaoianus (Theropoda: Dromaeosauridae) from Northeastern China." American Museum Novitates, 3381: 44pp.al.2002.pdf - Perle A. Norell M.A. and Clark J. 1999. A new maniraptoran theropod - Achillobator giganticus (Dromaeosauridae) - from the Upper Cretaceous of Burkhant, Mongolia. Contributions of the Mongolian-American Paleontological Project, 101: 1–105. - Britt, Chure, Stadtman, Madsen, Scheetz and Burge 2001. New osteological data and the affinities of Utahraptor from the Cedar Mountain Fm. (Early Cretaceous) of Utah. Journal of Vertebrate Paleontology, 21(3): 36A. - Naish, D. Hutt, and Martill D.M. 2001. Saurischian dinosaurs: theropods. in Martill D.M. and Naish D. (eds). Dinosaurs of the Isle of Wight. The Palaeontological Association, Field Guides to Fossils. 10, 242-309. - Turner, A.S.; Hwang S.H.; and Norell M.A. (2007). "A small derived theropod from Öösh, Early Cretaceous, Baykhangor Mongolia" (PDF). American Museum Novitates 3557: 1–27. doi:10.1206/0003-0082(2007)3557[1:ASDTFS]2.0.CO;2. http://digitallibrary.amnh.org/dspace/bitstream/2246/5845/1/N3557.pdf. Retrieved 2007-03-29. - Xing, Xu; et al. (2003). "Four-winged dinosaurs from China". Nature 421 (6921): 335–340. doi:10.1038/nature01342. PMID 12540892. - Xu X. Wang X.-L. and Wu X.-C. (1999). "A dromaeosaur dinosaur with a filamentous integument from the Yixian Formation of China". 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["Avian" features in the morphology of predatory dinosaurs]." Transactions of the Joint Soviet Mongolian Paleontological Expedition 24: 96–103. [Original article in Russian.] Translated by W. Robert Welsh, copy provided by Kenneth Carpenter and converted by Matthew Carrano. PDF fulltext - Z., Csiki; Vremir, M.; Brusatte, S. L.; and Norell, M. A. (in press). "An aberrant island-dwelling theropod dinosaur from the Late Cretaceous of Romania". Proceedings of the National Academy of Sciences of the United States of America preprint (35): 15357–61. doi:10.1073/pnas.1006970107. PMC 2932599. PMID 20805514. Supporting Information - Sweetman, S.C. (2004). "The first record of velociraptorine dinosaurs (Saurischia, Theropoda) from the Wealden (Early Cretaceous, Barremian) of southern England." Cretaceous Research, 25(3): 353–364. doi:10.1016/j.cretres.2004.01.004 - Bakker, Robert T. (1995). Raptor Red. New York: Bantam Books. pp. 4. ISBN 0-553-57561-9. Other websites [change] |Wikispecies has information on: Dromaeosauridae.| |Wikimedia Commons has media related to: Dromaeosauridae| - Dromaeosauridae at DinoData. - The Dromaeosauridae: The Raptors!, from the University of California Berkeley Museum of Paleontology. - Dromaeosauridae, by Justin Tweet from Thescelosaurus. - Dinosaurs – Complete and free online edition of the book "Dinosaurs" as written by W. D. Matthew (cited in this article with authorship of the family Dromaeosauridae), and former Curator of Vertebrate Paleontology at the American Museum of Natural History in New York; Originally published in 1915 - Dromaeosauridae, Dinosaur-world reference with in-depth description and pictures of many dromaeosauridae dinosaurs

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|This is the print version of Geometry You won't see this message or any elements not part of the book's content when you print or preview this page. Part I- Euclidean Geometry Chapter 1: Points, Lines, Line Segments and Rays Points and lines are two of the most fundamental concepts in Geometry, but they are also the most difficult to define. We can describe intuitively their characteristics, but there is no set definition for them: they, along with the plane, are the undefined terms of geometry. All other geometric definitions and concepts are built on the undefined ideas of the point, line and plane. Nevertheless, we shall try to define them. A point is an exact location in space. Points are dimensionless. That is, a point has no width, length, or height. We locate points relative to some arbitrary standard point, often called the "origin". Many physical objects suggest the idea of a point. Examples include the tip of a pencil, the corner of a cube, or a dot on a sheet of paper. As for a line segment, we specify a line with two points. Starting with the corresponding line segment, we find other line segments that share at least two points with the original line segment. In this way we extend the original line segment indefinitely. The set of all possible line segments findable in this way constitutes a line. A line extends indefinitely in a single dimension. Its length, having no limit, is infinite. Like the line segments that constitute it, it has no width or height. You may specify a line by specifying any two points within the line. For any two points, only one line passes through both points. On the other hand, an unlimited number of lines pass through any single point. We construct a ray similarly to the way we constructed a line, but we extend the line segment beyond only one of the original two points. A ray extends indefinitely in one direction, but ends at a single point in the other direction. That point is called the end-point of the ray. Note that a line segment has two end-points, a ray one, and a line none. A point exists in zero dimensions. A line exists in one dimension, and we specify a line with two points. A plane exists in two dimensions. We specify a plane with three points. Any two of the points specify a line. All possible lines that pass through the third point and any point in the line make up a plane. In more obvious language, a plane is a flat surface that extends indefinitely in its two dimensions, length and width. A plane has no height. Space exists in three dimensions. Space is made up of all possible planes, lines, and points. It extends indefinitely in all directions. Mathematics can extend space beyond the three dimensions of length, width, and height. We then refer to "normal" space as 3-dimensional space. A 4-dimensional space consists of an infinite number of 3-dimensional spaces. Etc. [How we label and reference points, lines, and planes.] Chapter 2: Angles An angle is the union of two rays with a common endpoint, called the vertex. The angles formed by vertical and horizontal lines are called right angles; lines, segments, or rays that intersect in right angles are said to be perpendicular. Angles, for our purposes, can be measured in either degrees (from 0 to 360) or radians (from 0 to ). Angles length can be determined by measuring along the arc they map out on a circle. In radians we consider the length of the arc of the circle mapped out by the angle. Since the circumference of a circle is , a right angle is radians. In degrees, the circle is 360 degrees, and so a right angle would be 90 degrees. Angles are named in several ways. - By naming the vertex of the angle (only if there is only one angle formed at that vertex; the name must be non-ambiguous) - By naming a point on each side of the angle with the vertex in between. - By placing a small number on the interior of the angle near the vertex. Classification of Angles by Degree Measure - an angle is said to be acute if it measures between 0 and 90 degrees, exclusive. - an angle is said to be right if it measures 90 degrees. - notice the small box placed in the corner of a right angle, unless the box is present it is not assumed the angle is 90 degrees. - all right angles are congruent - an angle is said to be obtuse if it measures between 90 and 180 degrees, exclusive. Special Pairs of Angles - adjacent angles - adjacent angles are angles with a common vertex and a common side. - adjacent angles have no interior points in common. - complementary angles - complementary angles are two angles whose sum is 90 degrees. - complementary angles may or may not be adjacent. - if two complementary angles are adjacent, then their exterior sides are perpendicular. - supplementary angles - two angles are said to be supplementary if their sum is 180 degrees. - supplementary angles need not be adjacent. - if supplementary angles are adjacent, then the sides they do not share form a line. - linear pair - if a pair of angles is both adjacent and supplementary, they are said to form a linear pair. - vertical angles - angles with a common vertex whose sides form opposite rays are called vertical angles. - vertical angles are congruent. Side-Side-Side (SSS) (Postulate 12) If three sides of one triangle are congruent to three sides of a second triangle, then the two triangles are congruent. Side-Angle-Side (SAS) (Postulate 13) If two sides and the included angle of a second triangle, then the two triangles are congruent. If two angles and the included side of one triangle are congruent to two angles and the included side of a second triangle, then two triangles are congruent. If two angles and a non-included side of one triangle are congruent to two angles and the corresponding non-included side of a second triangle, then the two triangles are congruent. NO - Angle-Side-Side (ASS) The "ASS" postulate does not work, unlike the other ones. A way that students can remember this is that "ass" is not a nice word, so we don't use it in geometry (since it does not work). There are two approaches to furthering knowledge: reasoning from known ideas and synthesizing observations. In inductive reasoning you observe the world, and attempt to explain based on your observations. You start with no prior assumptions. Deductive reasoning consists of logical assertions from known facts. What you need to know Before one can start to understand logic, and thereby begin to prove geometric theorems, one must first know a few vocabulary words and symbols. Conditional: a conditional is something which states that one statement implies another. A conditional contains two parts: the condition and the conclusion, where the former implies the latter. A conditional is always in the form "If statement 1, then statement 2." In most mathematical notation, a conditional is often written in the form p ⇒ q, which is read as "If p, then q" where p and q are statements. Converse: the converse of a logical statement is when the conclusion becomes the condition and vice versa; i.e., p ⇒ q becomes q ⇒ p. For example, the converse of the statement "If someone is a woman, then they are a human" would be "If someone is a human, then they are a woman." The converse of a conditional does not necessarily have the same truth value as the original, though it sometimes does, as will become apparent later. AND: And is a logical operator which is true only when both statements are true. For example, the statement "Diamond is the hardest substance known to man AND a diamond is a metal" is false. While the former statement is true, the latter is not. However, the statement "Diamond is the hardest substance known to man AND diamonds are made of carbon" would be true, because both parts are true. OR: If two statements are joined together by "or," then the truth of the "or" statement is dependant upon whether one or both of the statements from which it is composed is true. For example, the statement "Tuesday is the day after Monday OR Thursday is the day after Saturday" would have a truth value of "true," because even though the latter statement is false, the former is true. NOT: If a statement is preceded by "NOT," then it is evaluating the opposite truth value of that statement. The symbol for "NOT" is For example, if the statement p is "Elvis is dead," then ¬p would be "Elvis is not dead." The concept of "NOT" can cause some confusion when it relates to statements which contain the word "all." For example, if r is "¬". "All men have hair," then ¬r would be "All men do not have hair" or "No men have hair." Do not confuse this with "Not all men have hair" or "Some men have hair." The "NOT" should apply to the verb in the statement: in this case, "have." ¬p can also be written as NOT p or ~p. NOT p may also be referred to as the "negation of p." Inverse: The inverse of a conditional says that the negation of the condition implies the negation of the conclusion. For example, the inverse of p ⇒ q is ¬p ⇒ ¬q. Like a converse, an inverse does not necessarily have the same truth value as the original conditional. Biconditional: A biconditional is conditional where the condition and the conclusion imply one another. A biconditional starts with the words "if and only if." For example, "If and only if p, then q" means both that p implies q and that q implies p. Premise: A premise is a statement whose truth value is known initially. For example, if one were to say "If today is Thursday, then the cafeteria will serve burritos," and one knew that what day it was, then the premise would be "Today is Thursday" or "Today is not Thursday." ⇒: The symbol which denotes a conditional. p ⇒ q is read as "if p, then q." Iff: Iff is a shortened form of "if and only if." It is read as "if and only if." ⇔: The symbol which denotes a biconditonal. p ⇔ q is read as "If and only if p, then q." ∴: The symbol for "therefore." p ∴ q means that one knows that p is true (p is true is the premise), and has logically concluded that q must also be true. ∧: The symbol for "and." ∨: The symbol for "or." There are a few forms of deductive logic. One of the most common deductive logical arguments is modus ponens, which states that: - p ⇒ q - p ∴ q - (If p, then q) - (p, therefore q) An example of modus ponens: - If I stub my toe, then I will be in pain. - I stub my toe. - Therefore, I am in pain. Another form of deductive logic is modus tollens, which states the following. - p ⇒ q - ¬q ∴ ¬p - (If p, then q) - (not q, therefore not p) Modus tollens is just as valid a form of logic as modus ponens. The following is an example which uses modus tollens. - If today is Thursday, then the cafeteria will be serving burritos. - The cafeteria is not serving burritos, therefore today is not Thursday. Another form of deductive logic is known as the If-Then Transitive Property. Simply put, it means that there can be chains of logic where one thing implies another thing. The If-Then Transitive Property states: - p ⇒ q - (q ⇒ r) ∴ (p ⇒ r) - (If p, then q) - ((If q, then r), therefore (if p, then r)) For example, consider the following chain of if-then statements. - If today is Thursday, then the cafeteria will be serving burritos. - If the cafeteria will be serving burritos, then I will be happy. - Therefore, if today is Thursday, then I will be happy. Inductive reasoning is a logical argument which does not definitely prove a statement, but rather assumes it. Inductive reasoning is used often in life. Polling is an example of the use of inductive reasoning. If one were to poll one thousand people, and 300 of those people selected choice A, then one would infer that 30% of any population might also select choice A. This would be using inductive logic, because it does not definitively prove that 30% of any population would select choice A. Because of this factor of uncertainty, inductive reasoning should be avoided when possible when attempting to prove geometric properties. Truth tables are a way that one can display all the possibilities that a logical system may have when given certain premises. The following is a truth table with two premises (p and q), which shows the truth value of some basic logical statements. (NOTE: T = true; F = false) |p||q||¬p||¬q||p ⇒ q||p ⇔ q||p ∧ q||p ∨ q| Unlike science which has theories, mathematics has a definite notion of proof. Mathematics applies deductive reasoning to create a series of logical statements which show that one thing implies another. Consider a triangle, which we define as a shape with three vertices joined by three lines. We know that we can arbitrarily pick some point on a page, and make that into a vertex. We repeat that process and pick a second point. Using a ruler, we can connect these two points. We now make a third point, and using the ruler connect it to each of the other points. We have constructed a triangle. In mathematics we formalize this process into axioms, and carefully lay out the sequence of statements to show what follows. All definitions are clearly defined. In modern mathematics, we are always working within some system where various axioms hold. The most common form of explicit proof in highschool geometry is a two column proof consists of five parts: the given, the proposition, the statement column, the reason column, and the diagram (if one is given). Example of a Two-Column Proof Now, suppose a problem tells you to solve for , showing all steps made to get to the answer. A proof shows how this is done: Prove: x = 1 |Property of subtraction| We use "Given" as the first reason, because it is "given" to us in the problem. Written proofs (also known as informal proofs, paragraph proofs, or 'plans for proof') are written in paragraph form. Other than this formatting difference, they are similar to two-column proofs. Sometimes it is helpful to start with a written proof, before formalizing the proof in two-column form. If you're having trouble putting your proof into two column form, try "talking it out" in a written proof first. Example of a Written Proof We are given that x + 1 = 2, so if we subtract one from each side of the equation (x + 1 - 1 = 2 - 1), then we can see that x = 1 by the definition of subtraction. A flowchart proof or more simply a flow proof is a graphical representation of a two-column proof. Each set of statement and reasons are recorded in a box and then arrows are drawn from one step to another. This method shows how different ideas come together to formulate the proof. Postulates in geometry are very similar to axioms, self-evident truths, and beliefs in logic, political philosophy and personal decision-making. The five postulates of Euclidean Geometry define the basic rules governing the creation and extension of geometric figures with ruler and compass. Together with the five axioms (or "common notions") and twenty-three definitions at the beginning of Euclid's Elements, they form the basis for the extensive proofs given in this masterful compilation of ancient Greek geometric knowledge. They are as follows: - A straight line may be drawn from any given point to any other. - A straight line may be extended to any finite length. - A circle may be described with any given point as its center and any distance as its radius. - All right angles are congruent. - If a straight line intersects two other straight lines, and so makes the two interior angles on one side of it together less than two right angles, then the other straight lines will meet at a point if extended far enough on the side on which the angles are less than two right angles. Postulate 5, the so-called Parallel Postulate was the source of much annoyance, probably even to Euclid, for being so relatively prolix. Mathematicians have a peculiar sense of aesthetics that values simplicity arising from simplicity, with the long complicated proofs, equations and calculations needed for rigorous certainty done behind the scenes, and to have such a long sentence amidst such other straightforward, intuitive statements seems awkward. As a result, many mathematicians over the centuries have tried to prove the results of the Elements without using the Parallel Postulate, but to no avail. However, in the past two centuries, assorted non-Euclidean geometries have been derived based on using the first four Euclidean postulates together with various negations of the fifth. Chapter 7. Vertical Angles Vertical angles are a pair of angles with a common vertex whose sides form opposite rays. An extensively useful fact about vertical angles is that they are congruent. Aside from saying that any pair of vertical angles "obviously" have the same measure by inspection, we can prove this fact with some simple algebra and an observation about supplementary angles. Let two lines intersect at a point, and angles A1 and A2 be a pair of vertical angles thus formed. At the point of intersection, two other angles are also formed, and we'll call either one of them B1 without loss of generality. Since B1 and A1 are supplementary, we can say that the measure of B1 plus the measure of A1 is 180. Similarly, the measure of B1 plus the measure of A2 is 180. Thus the measure of A1 plus the measure of B1 equals the measure of A2 plus the measure of B1, by substitution. Then by subracting the measure of B1 from each side of this equality, we have that the measure of A1 equals the measure of A2. Parallel Lines in a Plane Two coplanar lines are said to be parallel if they never intersect. For any given point on the first line, its distance to the second line is equal to the distance between any other point on the first line and the second line. The common notation for parallel lines is "||" (a double pipe); it is not unusual to see "//" as well. If line m is parallel to line n, we write "m || n". Lines in a plane either coincide, intersect in a point, or are parallel. Controversies surrounding the Parallel Postulate lead to the development of non-Euclidean geometries. Parallel Lines and Special Pairs of Angles When two (or more) parallel lines are cut by a transversal, the following angle relationships hold: - corresponding angles are congruent - alternate exterior angles are congruent - same-side interior angles are supplementary Theorems Involving Parallel Lines - If a line in a plane is perpendicular to one of two parallel lines, it is perpendicular to the other line as well. - If a line in a plane is parallel to one of two parallel lines, it is parallel to both parallel lines. - If three or more parallel lines are intersected by two or more transversals, then they divide the transversals proportionally. Congruent shapes are the same size with corresponding lengths and angles equal. In other words, they are exactly the same size and shape. They will fit on top of each other perfectly. Therefore if you know the size and shape of one you know the size and shape of the others. For example: Each of the above shapes is congruent to each other. The only difference is in their orientation, or the way they are rotated. If you traced them onto paper and cut them out, you could see that they fit over each other exactly. Having done this, right away we can see that, though the angles correspond in size and position, the sides do not. Therefore it is proved the triangles are not congruent. Similar shapes are like congruent shapes in that they must be the same shape, but they don't have to be the same size. Their corresponding angles are congruent and their corresponding sides are in proportion. Methods of Determining Congruence Two triangles are congruent if: - each pair of corresponding sides is congruent - two pairs of corresponding angles are congruent and a pair of corresponding sides are congruent - two pairs of corresponding sides and the angles included between them are congruent Tips for Proofs Commonly used prerequisite knowledge in determining the congruence of two triangles includes: - by the reflexive property, a segment is congruent to itself - vertical angles are congruent - when parallel lines are cut by a transversal corresponding angles are congruent - when parallel lines are cut by a transversal alternate interior angles are congruent - midpoints and bisectors divide segments and angles into two congruent parts For two triangles to be similar, all 3 corresponding angles must be congruent, and all three sides must be proportionally equal. Two triangles are similar if... - Two angles of each triangle are congruent. - The acute angle of a right triangle is congruent to the acute angle of another right triangle. - The two triangles are congruent. Note here that congruency implies similarity. A quadrilateral is a polygon that has four sides. Special Types of Quadrilaterals - A parallelogram is a quadrilateral having two pairs of parallel sides. - A square, a rhombus, and a rectangle are all examples of parallelograms. - A rhombus is a quadrilateral of which all four sides are the same length. - A rectangle is a parallelogram of which all four angles are 90 degrees. - A square is a quadrilateral of which all four sides are of the same length, and all four angles are 90 degrees. - A square is a rectangle, a rhombus, and a parallelogram. - A trapezoid is a quadrilateral which has two parallel sides (U.S.) - U.S. usage: A trapezium is a quadrilateral which has no parallel sides. - U.K usage: A trapezium is a quadrilateral with two parallel sides (same as US trapezoid definition). - A kite is an quadrilateral with two pairs of congruent adjacent sides. One of the most important properties used in proofs is that the sum of the angles of the quadrilateral is always 360 degrees. This can easily be proven too: If you draw a random quadrilateral, and one of its diagonals, you'll split it up into two triangles. Given that the sum of the angles of a triangle is 180 degrees, you can sum them up, and it'll give 360 degrees. A parallelogram is a geometric figure with two pairs of parallel sides. Parallelograms are a special type of quadrilateral. The opposite sides are equal in length and the opposite angles are also equal. The area is equal to the product of any side and the distance between that side and the line containing the opposite side. Properties of Parallelograms The following properties are common to all parallelograms (parallelogram, rhombus, rectangle, square) - both pairs of opposite sides are parallel - both pairs of opposite sides are congruent - both pairs of opposite angles are congruent - the diagonals bisect each other - A rhombus is a parallelogram with four congruent sides. - The diagonals of a rhombus are perpendicular. - Each diagonal of a rhombus bisects two angles the rhombus. - A rhombus may or may not be a square. - A square is a parallelogram with four right angles and four congruent sides. - A square is both a rectangle and a rhombus and inherits all of their properties. A Trapezoid (American English) or Trapezium (British English) is a quadrilateral that has two parallel sides and two non parallel sides. Some properties of trapezoids: - The interior angles sum to 360° as in any quadrilateral. - The parallel sides are unequal. - Each of the parallel sides is called a base (b) of the trapezoid. The two angles that join one base are called 'base angles'. - If the two non-parallel sides are equal, the trapezoid is called an isosceles trapezoid. - In an isosceles trapezoid, each pair of base angles are equal. - If one pair of base angles of a trapezoid are equal, the trapezoid is isosceles. - A line segment connecting the midpoints of the non-parallel sides is called the median (m) of the trapeziod. - The median of a trapezoid is equal to one half the sum of the bases (called b1 and b2). - A line segment perpendicular to the bases is called an altitude (h) of the trapezoid. The area (A) of a trapezoid is equal to the product of an altitude and the median. Recall though that the median is half of the sum of the bases. Substituting for m, we get: A circle is a set of all points in a plane that are equidistant from a single point; that single point is called the centre of the circle and the distance between any point on circle and the centre is called radius of the circle. a chord is an internal segment of a circle that has both of its endpoints on the circumference of the circle. - the diameter of a circle is the largest chord possible a secant of a circle is any line that intersects a circle in two places. - a secant contains any chord of the circle a tangent to a circle is a line that intersects a circle in exactly one point, called the point of tangency. - at the point of tangency the tangent line and the radius of the circle are perpendicular Chapter 16. Circles/Arcs An arc is a segment of the perimeter of a given circle. The measure of an arc is measured as an angle, this could be in radians or degrees (more on radians later). The exact measure of the arc is determined by the measure of the angle formed when a line is drawn from the center of the circle to each end point. As an example the circle below has an arc cut out of it with a measure of 30 degrees. As I mentioned before an arc can be measured in degrees or radians. A radian is merely a different method for measuring an angle. If we take a unit circle (which has a radius of 1 unit), then if we take an arc with the length equal to 1 unit, and draw line from each endpoint to the center of the circle the angle formed is equal to 1 radian. this concept is displayed below, in this circle an arc has been cut off by an angle of 1 radian, and therefore the length of the arc is equal to because the radius is 1. From this definition we can say that on the unit circle a single radian is equal to radians because the perimeter of a unit circle is equal to . Another useful property of this definition that will be extremely useful to anyone who studies arcs is that the length of an arc is equal to its measure in radians multiplied by the radius of the circle. Converting to and from radians is a fairly simple process. 2 facts are required to do so, first a circle is equal to 360 degrees, and it is also equal to . using these 2 facts we can form the following formula: , thus 1 degree is equal to radians. From here we can simply multiply by the number of degrees to convert to radians. for example if we have 20 degrees and want to convert to radians then we proceed as follows: The same sort of argument can be used to show the formula for getting 1 radian. , thus 1 radian is equal to A tangent is a line in the same plane as a given circle that meets that circle in exactly one point. That point is called the point of tangency. A tangent cannot pass through a circle; if it does, it is classified as a chord. A secant is a line containing a chord. A common tangent is a line tangent to two circles in the same plane. If the tangent does not intersect the line containing and connecting the centers of the circles, it is an external tangent. If it does, it is an internal tangent. Two circles are tangent to one another if in a plane they intersect the same tangent in the same point. Sector of a circle A sector of a circle can be thought of as a pie piece. In the picture below, a sector of the circle is shaded yellow. To find the area of a sector, find the area of the whole circle and then multiply by the angle of the sector over 360 degrees. A more intuitive approach can be used when the sector is half the circle. In this case the area of the sector would just be the area of the circle divided by 2. - See Angle Addition Property of Equality For any real numbers a, b, and c, if a = b, then a + c = b + c. A figure is an angle if and only if it is composed of two rays which share a common endpoint. Each of these rays (or segments, as the case may be) is known as a side of the angle (For example, in the illustration at right), and the common point is known as the angle's vertex (point B in the illustration). Angles are measured by the difference of their slopes. The units for angle measure are radians and degrees. Angles may be classified by their degree measure. - Acute Angle: an angle is an acute angle if and only if it has a measure of less than 90° - Right Angle: an angle is an right angle if and only if it has a measure of exactly 90° - Obtuse Angle: an angle is an obtuse angle if and only if it has a measure of greater than 90° Angle Addition Postulate If P is in the interior of an angle , then Center of a circle Point P is the center of circle C if and only if all points in circle C are equidistant from point P and point P is contained in the same plane as circle C. A collection of points is said to be a circle with a center at point P and a radius of some distance r if and only if it is the collection of all points which are a distance of r away from point P and are contained by a plane which contain point P. A polygon is said to be concave if and only if it contains at least one interior angle with a measure greater than 180° exclusively and less than 360° exclusively. Two angles formed by a transversal intersecting with two lines are corresponding angles if and only if one is on the inside of the two lines, the other is on the outside of the two lines, and both are on the same side of the transversal. Corresponding Angles Postulate If two lines cut by a transversal are parallel, then their corresponding angles are congruent. Corresponding Parts of Congruent Triangles are Congruent Postulate The Corresponding Parts of Congruent Triangles are Congruent Postulate (CPCTC) states: - If ∆ABC ≅ ∆XYZ, then all parts of ∆ABC are congruent to their corresponding parts in ∆XYZ. For example: - ∠ABC ≅ ∠XYZ - ∠BCA ≅ ∠YZX - ∠CAB ≅ ∠ZXY CPCTC also applies to all other parts of the triangles, such as a triangle's altitude, median, circumcenter, et al. A line segment is the diameter of a circle if and only if it is a chord of the circle which contains the circle's center. - See Circle and if they cross they are congruent A collection of points is a line if and only if the collection of points is perfectly straight (aligned), is infinitely long, and is infinitely thin. Between any two points on a line, there exists an infinite number of points which are also contained by the line. Lines are usually written by two points in the line, such as line AB, or A collection of points is a line segment if and only if it is perfectly straight, is infinitely thin, and has a finite length. A line segment is measured by the shortest distance between the two extreme points on the line segment, known as endpoints. Between any two points on a line segment, there exists an infinite number of points which are also contained by the line segment. Two lines or line segments are said to be parallel if and only if the lines are contained by the same plane and have no points in common if continued infinitely. Two planes are said to be parallel if and only if the planes have no points in common when continued infinitely. Two lines that intersect at a 90° angle. Given a line, and a point P not in line , then there is one and only one line that goes through point P perpendicular to An object is a plane if and only if it is a two-dimensional object which has no thickness or curvature and continues infinitely. A plane can be defined by three points. A plane may be considered to be analogous to a piece of paper. A point is a zero-dimensional mathematical object representing a location in one or more dimensions. A point has no size; it has only location. A polygon is a closed plane figure composed of at least 3 straight lines. Each side has to intersect another side at their respective endpoints, and that the lines intersecting are not collinear. The radius of a circle is the distance between any given point on the circle and the circle's center. - See Circle A ray is a straight collection of points which continues infinitely in one direction. The point at which the ray stops is known as the ray's endpoint. Between any two points on a ray, there exists an infinite number of points which are also contained by the ray. The points on a line can be matched one to one with the real numbers. The real number that corresponds to a point is the point's coordinate. The distance between two points is the absolute value of the difference between the two coordinates of the two points. Geometry/Synthetic versus analytic geometry - Two and Three-Dimensional Geometry and Other Geometric Figures Perimeter and Arclength Perimeter of Circle The circles perimeter can be calculated using the following formula where and the radius of the circle. Perimeter of Polygons The perimeter of a polygon with number of sides abbreviated can be caculated using the following formula Arclength of Circles The arclength of a given circle with radius can be calculated using where is the angle given in radians. Arclength of Curves If a curve in have a parameter form for , then the arclength can be calculated using the following fomula Derivation of formula can be found using differential geometry on infinitely small triangles. Area of Circles The method for finding the area of a circle is Where π is a constant roughly equal to 3.14159265358978 and r is the radius of the circle; a line drawn from any point on the circle to its center. Area of Triangles Three ways of calculating the area inside of a triangle are mentioned here. If one of the sides of the triangle is chosen as a base, then a height for the triangle and that particular base can be defined. The height is a line segment perpendicular to the base or the line formed by extending the base and the endpoints of the height are the corner point not on the base and a point on the base or line extending the base. Let B = the length of the side chosen as the base. Let h = the distance between the endpoints of the height segment which is perpendicular to the base. Then the area of the triangle is given by: This method of calculating the area is good if the value of a base and its corresponding height in the triangle is easily determined. This is particularly true if the triangle is a right triangle, and the lengths of the two sides sharing the 90o angle can be determined. - , also known as Heron's Formula If the lengths of all three sides of a triangle are known, Hero's formula may be used to calculate the area of the triangle. First, the semiperimeter, s, must be calculated by dividing the sum of the lengths of all three sides by 2. For a triangle having side lengths a, b, and c : Then the triangle's area is given by: If the triangle is needle shaped, that is, one of the sides is very much shorter than the other two then it can be difficult to compute the area because the precision needed is greater than that available in the calculator or computer that is used. In otherwords Heron's formula is numerically unstable. Another formula that is much more stable is: where , , and have been sorted so that . In a triangle with sides length a, b, and c and angles A, B, and C opposite them, This formula is true because in the formula . It is useful because you don't need to find the height from an angle in a separate step, and is also used to prove the law of sines (divide all terms in the above equation by a*b*c and you'll get it directly!) Area of Rectangles The area calculation of a rectangle is simple and easy to understand. One of the sides is chosen as the base, with a length b. An adjacent side is then the height, with a length h, because in a rectangle the adjacent sides are perpendicular to the side chosen as the base. The rectangle's area is given by: Sometimes, the baselength may be referred to as the length of the rectangle, l, and the height as the width of the rectangle, w. Then the area formula becomes: Regardless of the labels used for the sides, it is apparent that the two formulas are equivalent. Of course, the area of a square with sides having length s would be: Area of Parallelograms The area of a parallelogram can be determined using the equation for the area of a rectangle. The formula is: A is the area of a parallelogram. b is the base. h is the height. The height is a perpendicular line segment that connects one of the vertices to its opposite side (the base). Area of Rhombus Remember in a rombus all sides are equal in length. and represent the diagonals. Area of Trapezoids The area of a trapezoid is derived from taking the arithmetic mean of its two parallel sides to form a rectangle of equal area. Where and are the lengths of the two parallel bases. Area of Kites The area of a kite is based on splitting the kite into four pieces by halving it along each diagonal and using these pieces to form a rectangle of equal area. Where a and b are the diagonals of the kite. Alternatively, the kite may be divided into two halves, each of which is a triangle, by the longer of its diagonals, a. The area of each triangle is thus Where b is the other (shorter) diagonal of the kite. And the total area of the kite (which is composed of two identical such triangles) is Which is the same as Areas of other Quadrilaterals The areas of other quadrilaterals are slightly more complex to calculate, but can still be found if the quadrilateral is well-defined. For example, a quadrilateral can be divided into two triangles, or some combination of triangles and rectangles. The areas of the constituent polygons can be found and added up with arithmetic. Volume is like area expanded out into 3 dimensions. Area deals with only 2 dimensions. For volume we have to consider another dimension. Area can be thought of as how much space some drawing takes up on a flat piece of paper. Volume can be thought of as how much space an object takes up. |Common equations for volume:| |A cube:||s = length of a side| |A rectangular prism:||l = length, w = width, h = height| |A cylinder (circular prism):||r = radius of circular face, h = height| |Any prism that has a constant cross sectional area along the height:||A = area of the base, h = height| |A sphere:||r = radius of sphere which is the integral of the Surface Area of a sphere |An ellipsoid:||a, b, c = semi-axes of ellipsoid| |A pyramid:||A = area of the base, h = height of pyramid| |A cone (circular-based pyramid):||r = radius of circle at base, h = distance from base to tip (The units of volume depend on the units of length - if the lengths are in meters, the volume will be in cubic meters, etc.) The volume of any solid whose cross sectional areas are all the same is equal to that cross sectional area times the distance the centroid(the center of gravity in a physical object) would travel through the solid. If two solids are contained between two parallel planes and every plane parallel to these two plane has equal cross sections through these two solids, then their volumes are equal. A Polygon is a two-dimensional figure, meaning all of the lines in the figure are contained within one plane. They are classified by the number of angles, which is also the number of sides. One key point to note is that a polygon must have at least three sides. Normally, three to ten sided figures are referred to by their names (below), while figures with eleven or more sides is an n-gon, where n is the number of sides. Hence a forty-sided polygon is called a 40-gon. A polygon with three angles and sides. A polygon with four angles and sides. A polygon with five angles and sides. A polygon with six angles and sides. A polygon with seven angles and sides. A polygon with eight angles and sides. A polygon with nine angles and sides. A polygon with ten angles and sides. For a list of n-gon names, go to and scroll to the bottom of the page. Polygons are also classified as convex or concave. A convex polygon has interior angles less than 180 degrees, thus all triangles are convex. If a polygon has at least one internal angle greater than 180 degrees, then it is concave. An easy way to tell if a polygon is concave is if one side can be extended and crosses the interior of the polygon. Concave polygons can be divided into several convex polygons by drawing diagonals. Regular polygons are polygons in which all sides and angles are congruent. A triangle is a type of polygon having three sides and, therefore, three angles. The triangle is a closed figure formed from three straight line segments joined at their ends. The points at the ends can be called the corners, angles, or vertices of the triangle. Since any given triangle lies completely within a plane, triangles are often treated as two-dimensional geometric figures. As such, a triangle has no volume and, because it is a two-dimensionally closed figure, the flat part of the plane inside the triangle has an area, typically referred to as the area of the triangle. Triangles are always convex polygons. A triangle must have at least some area, so all three corner points of a triangle cannot lie in the same line. The sum of the lengths of any two sides of a triangle is always greater than the length of the third side. The preceding statement is sometimes called the Triangle Inequality. Certain types of triangles Categorized by angle The sum of the interior angles in a triangle always equals 180o. This means that no more than one of the angles can be 90o or more. All three angles can all be less than 90oin the triangle; then it is called an acute triangle. One of the angles can be 90o and the other two less than 90o; then the triangle is called a right triangle. Finally, one of the angles can be more than 90o and the other two less; then the triangle is called an obtuse triangle. Categorized by sides If all three of the sides of a triangle are of different length, then the triangle is called a scalene triangle. If two of the sides of a triangle are of equal length, then it is called an isoceles triangle. In an isoceles triangle, the angle between the two equal sides can be more than, equal to, or less than 90o. The other two angles are both less than 90o. If all three sides of a triangle are of equal length, then it is called an equilateral triangle and all three of the interior angles must be 60o, making it equilangular. Because the interior angles are all equal, all equilateral triangles are also the three-sided variety of a regular polygon and they are all similar, but might not be congruent. However, polygons having four or more equal sides might not have equal interior angles, might not be regular polygons, and might not be similar or congruent. Of course, pairs of triangles which are not equilateral might be similar or congruent. Opposite corners and sides in triangles If one of the sides of a triangle is chosen, the interior angles of the corners at the side's endpoints can be called adjacent angles. The corner which is not one of these endpoints can be called the corner opposite to the side. The interior angle whose vertex is the opposite corner can be called the angle opposite to the side. Likewise, if a corner or its angle is chosen, then the two sides sharing an endpoint at that corner can be called adjacent sides. The side not having this corner as one of its two endpoints can be called the side opposite to the corner. The sides or their lengths of a triangle are typically labeled with lower case letters. The corners or their corresponding angles can be labeled with capital letters. The triangle as a whole can be labeled by a small triangle symbol and its corner points. In a triangle, the largest interior angle is opposite to longest side, and vice versa. Any triangle can be divided into two right triangles by taking the longest side as a base, and extending a line segment from the opposite corner to a point on the base such that it is perpendicular to the base. Such a line segment would be considered the height or altitude ( h ) for that particular base ( b ). The two right triangles resulting from this division would both share the height as one of its sides. The interior angles at the meeting of the height and base would be 90o for each new right triangle. For acute triangles, any of the three sides can act as the base and have a corresponding height. For more information on right triangles, see Right Triangles and Pythagorean Theorem. Area of Triangles If base and height of a triangle are known, then the area of the triangle can be calculated by the formula: ( is the symbol for area) Ways of calculating the area inside of a triangle are further discussed under Area. The centroid is constructed by drawing all the medians of the triangle. All three medians intersect at the same point: this crossing point is the centroid. Centroids are always inside a triangle. They are also the centre of gravity of the triangle. The three angle bisectors of the triangle intersect at a single point, called the incentre. Incentres are always inside the triangle. The three sides are equidistant from the incentre. The incentre is also the centre of the inscribed circle (incircle) of a triangle, or the interior circle which touches all three sides of the triangle. The circumcentre is the intersection of all three perpendicular bisectors. Unlike the incentre, it is outside the triangle if the triangle is obtuse. Acute triangles always have circumcentres inside, while the circumcentre of a right triangle is the midpoint of the hypotenuse. The vertices of the triangle are equidistant from the circumcentre. The circumcentre is so called because it is the centre of the circumcircle, or the exterior circle which touches all three vertices of the triangle. The orthocentre is the crossing point of the three altitudes. It is always inside acute triangles, outside obtuse triangles, and on the right vertex of the right-angled triangle. Please note that the centres of an equilateral triangle are always the same point. Right Triangles and Pythagorean Theorem Right triangles are triangles in which one of the interior angles is 90o. A 90o angle is called a right angle. Right triangles are sometimes called right-angled triangles. The other two interior angles are complementary, i.e. their sum equals 90o. Right triangles have special properties which make it easier to conceptualize and calculate their parameters in many cases. The side opposite of the right angle is called the hypotenuse. The sides adjacent to the right angle are the legs. When using the Pythagorean Theorem, the hypotenuse or its length is often labeled with a lower case c. The legs (or their lengths) are often labeled a and b. Either of the legs can be considered a base and the other leg would be considered the height (or altitude), because the right angle automatically makes them perpendicular. If the lengths of both the legs are known, then by setting one of these sides as the base ( b ) and the other as the height ( h ), the area of the right triangle is very easy to calculate using this formula: This is intuitively logical because another congruent right triangle can be placed against it so that the hypotenuses are the same line segment, forming a rectangle with sides having length b and width h. The area of the rectangle is b × h, so either one of the congruent right triangles forming it has an area equal to half of that rectangle. Right triangles can be neither equilateral, acute, nor obtuse triangles. Isosceles right triangles have two 45° angles as well as the 90° angle. All isosceles right triangles are similar since corresponding angles in isosceles right triangles are equal. If another triangle can be divided into two right triangles (see Triangle), then the area of the triangle may be able to be determined from the sum of the two constituent right triangles. Also the Pythagorean theorem can be used for non right triangles. a2+b2=c2-2c For history regarding the Pythagorean Theorem, see Pythagorean theorem. The Pythagorean Theorem states that: - In a right triangle, the square of the length of the hypotenuse is equal to the sum of the squares of the lengths of the other two sides. Let's take a right triangle as shown here and set c equal to the length of the hypotenuse and set a and b each equal to the lengths of the other two sides. Then the Pythagorean Theorem can be stated as this equation: Using the Pythagorean Theorem, if the lengths of any two of the sides of a right triangle are known and it is known which side is the hypotenuse, then the length of the third side can be determined from the formula. Sine, Cosine, and Tangent for Right Triangles Sine, Cosine, and Tangent are all functions of an angle, which are useful in right triangle calculations. For an angle designated as θ, the sine function is abbreviated as sin θ, the cosine function is abbreviated as cos θ, and the tangent function is abbreviated as tan θ. For any angle θ, sin θ, cos θ, and tan θ are each single determined values and if θ is a known value, sin θ, cos θ, and tan θ can be looked up in a table or found with a calculator. There is a table listing these function values at the end of this section. For an angle between listed values, the sine, cosine, or tangent of that angle can be estimated from the values in the table. Conversely, if a number is known to be the sine, cosine, or tangent of a angle, then such tables could be used in reverse to find (or estimate) the value of a corresponding angle. These three functions are related to right triangles in the following ways: In a right triangle, - the sine of a non-right angle equals the length of the leg opposite that angle divided by the length of the hypotenuse. - the cosine of a non-right angle equals the length of the leg adjacent to it divided by the length of the hypotenuse. - the tangent of a non-right angle equals the length of the leg opposite that angle divided by the length of the leg adjacent to it. For any value of θ where cos θ ≠ 0, If one considers the diagram representing a right triangle with the two non-right angles θ1and θ2, and the side lengths a,b,c as shown here: For the functions of angle θ1: Analogously, for the functions of angle θ2: Table of sine, cosine, and tangent for angles θ from 0 to 90° |θ in degrees||θ in radians||sin θ||cos θ||tan θ| General rules for important angles: Polyominoes are shapes made from connecting unit squares together, though certain connections are not allowed. A domino is the shape made from attaching unit squares so that they share one full edge. The term polyomino is based on the word domino. There is only one possible domino. Tromino↑Jump back a section A polymino made from four squares is called a tetromino. There are five possible combinations and two reflections: A polymino made from five squares is called a pentomino. There are twelve possible pentominoes, excluding mirror images and rotations. Ellipses are sometimes called ovals. Ellipses contain two foci. The sum of the distance from a point on the ellipse to one focus and that same point to the other focus is constant Area Shapes Extended into 3rd Dimension Geometry/Area Shapes Extended into 3rd Dimension Area Shapes Extended into 3rd Dimension Linearly to a Line or Point Geometry/Area Shapes Extended into 3rd Dimension Linearly to a Line or Point Ellipsoids and Spheres Geometry/Ellipsoids and Spheres Suppose you are an astronomer in America. You observe an exciting event (say, a supernova) in the sky and would like to tell your colleagues in Europe about it. Suppose the supernova appeared at your zenith. You can't tell astronomers in Europe to look at their zenith because their zenith points in a different direction. You might tell them which constellation to look in. This might not work, though, because it might be too hard to find the supernova by searching an entire constellation. The best solution would be to give them an exact position by using a coordinate system. On Earth, you can specify a location using latitude and longitude. This system works by measuring the angles separating the location from two great circles on Earth (namely, the equator and the prime meridian). Coordinate systems in the sky work in the same way. The equatorial coordinate system is the most commonly used. The equatorial system defines two coordinates: right ascension and declination, based on the axis of the Earth's rotation. The declination is the angle of an object north or south of the celestial equator. Declination on the celestial sphere corresponds to latitude on the Earth. The right ascension of an object is defined by the position of a point on the celestial sphere called the vernal equinox. The further an object is east of the vernal equinox, the greater its right ascension. A coordinate system is a system designed to establish positions with respect to given reference points. The coordinate system consists of one or more reference points, the styles of measurement (linear measurement or angular measurement) from those reference points, and the directions (or axes) in which those measurements will be taken. In astronomy, various coordinate systems are used to precisely define the locations of astronomical objects. Latitude and longitude are used to locate a certain position on the Earth's surface. The lines of latitude (horizontal) and the lines of longitude (vertical) make up an invisible grid over the earth. Lines of latitude are called parallels. Lines of longitude aren't completely straight (they run from the exact point of the north pole to the exact point of the south pole) so they are called meridians. 0 degrees latitude is the Earth's middle, called the equator. O degrees longitude was tricky because there really is no middle of the earth vertically. It was finally agreed that the observatory in Greenwich, U.K. would be 0 degrees longitude due to its significant roll in scientific discoveries and creating latitude and longitude. 0 degrees longitude is called the prime meridian. Latitude and longitude are measured in degrees. One degree is about 69 miles. There are sixty minutes (') in a degree and sixty seconds (") in a minute. These tiny units make GPS's (Global Positioning Systems) much more exact. There are a few main lines of latitude:the Arctic Circle, the Antarctic Circle, the Tropic of Cancer, and the Tropic of Capricorn. The Antarctic Circle is 66.5 degrees south of the equator and it marks the temperate zone from the Antarctic zone. The Arctic Circle is an exact mirror in the north. The Tropic of Cancer separates the tropics from the temperate zone. It is 23.5 degrees north of the equator. It is mirrored in the south by the Tropic of Capricorn. Horizontal coordinate system One of the simplest ways of placing a star on the night sky is the coordinate system based on altitude or azimuth, thus called the Alt-Az or horizontal coordinate system. The reference circles for this system are the horizon and the celestial meridian, both of which may be most easily graphed for a given location using the celestial sphere. In simplest terms, the altitude is the angle made from the position of the celestial object (e.g. star) to the point nearest it on the horizon. The azimuth is the angle from the northernmost point of the horizon (which is also its intersection with the celestial meridian) to the point on the horizon nearest the celestial object. Usually azimuth is measured eastwards from due north. So east has az=90°, south has az=180°, west has az=270° and north has az=360° (or 0°). An object's altitude and azimuth change as the earth rotates. Equatorial coordinate system The equatorial coordinate system is another system that uses two angles to place an object on the sky: right ascension and declination. Ecliptic coordinate system The ecliptic coordinate system is based on the ecliptic plane, i.e., the plane which contains our Sun and Earth's average orbit around it, which is tilted at 23°26' from the plane of Earth's equator. The great circle at which this plane intersects the celestial sphere is the ecliptic, and one of the coordinates used in the ecliptic coordinate system, the ecliptic latitude, describes how far an object is to ecliptic north or to ecliptic south of this circle. On this circle lies the point of the vernal equinox (also called the first point of Aries); ecliptic longitude is measured as the angle of an object relative to this point to ecliptic east. Ecliptic latitude is generally indicated by φ, whereas ecliptic longitude is usually indicated by λ. Galactic coordinate system As a member of the Milky Way Galaxy, we have a clear view of the Milky Way from Earth. Since we are inside the Milky Way, we don't see the galaxy's spiral arms, central bulge and so forth directly as we do for other galaxies. Instead, the Milky Way completely encircles us. We see the Milky Way as a band of faint starlight forming a ring around us on the celestial sphere. The disk of the galaxy forms this ring, and the bulge forms a bright patch in the ring. You can easily see the Milky Way's faint band from a dark, rural location. Our galaxy defines another useful coordinate system — the galactic coordinate system. This system works just like the others we've discussed. It also uses two coordinates to specify the position of an object on the celestial sphere. The galactic coordinate system first defines a galactic latitude, the angle an object makes with the galactic equator. The galactic equator has been selected to run through the center of the Milky Way's band. The second coordinate is galactic longitude, which is the angular separation of the object from the galaxy's "prime meridian," the great circle that passes through the Galactic center and the galactic poles. The galactic coordinate system is useful for describing an object's position with respect to the galaxy's center. For example, if an object has high galactic latitude, you might expect it to be less obstructed by interstellar dust. Transformations between coordinate systems One can use the principles of spherical trigonometry as applied to triangles on the celestial sphere to derive formulas for transforming coordinates in one system to those in another. These formulas generally rely on the spherical law of cosines, known also as the cosine rule for sides. By substituting various angles on the celestial sphere for the angles in the law of cosines and by thereafter applying basic trigonometric identities, most of the formulas necessary for coordinate transformations can be found. The law of cosines is stated thus: To transform from horizontal to equatorial coordinates, the relevant formulas are as follows: where RA is the right ascension, Dec is the declination, LST is the local sidereal time, Alt is the altitude, Az is the azimuth, and Lat is the observer's latitude. Using the same symbols and formulas, one can also derive formulas to transform from equatorial to horizontal coordinates: Transformation from equatorial to ecliptic coordinate systems can similarly be accomplished using the following formulas: where RA is the right ascension, Dec is the declination, φ is the ecliptic latitude, λ is the ecliptic longitude, and ε is the tilt of Earth's axis relative to the ecliptic plane. Again, using the same formulas and symbols, new formulas for transforming ecliptic to equatorial coordinate systems can be found: - Traditional Geometry: A topological space is a set X, and a collection of subsets of X, C such that both the empty set and X are contained in C and the union of any subcollection of sets in C and the intersection of any finite subcollection of sets in C are also contained within C. The sets in C are called open sets. Their complements relative to X are called closed sets. Given two topological spaces, X and Y, a map f from X to Y is continuous if for every open set U of Y, f−1(U) is an open set of X. Hyperbolic and Elliptic Geometry There are precisely three different classes of three-dimensional constant-curvature geometry: Euclidean, hyperbolic and elliptic geometry. The three geometries are all built on the same first four axioms, but each has a unique version of the fifth axiom, also known as the parallel postulate. The 1868 Essay on an Interpretation of Non-Euclidean Geometry by Eugenio Beltrami (1835 - 1900) proved the logical consistency of the two Non-Euclidean geometries, hyperbolic and elliptic. The Parallel Postulate The parallel postulate is as follows for the corresponding geometries. Euclidean geometry: Playfair's version: "Given a line l and a point P not on l, there exists a unique line m through P that is parallel to l." Euclid's version: "Suppose that a line l meets two other lines m and n so that the sum of the interior angles on one side of l is less than 180°. Then m and n intersect in a point on that side of l." These two versions are equivalent; though Playfair's may be easier to conceive, Euclid's is often useful for proofs. Hyperbolic geometry: Given an arbitrary infinite line l and any point P not on l, there exist two or more distinct lines which pass through P and are parallel to l. Elliptic geometry: Given an arbitrary infinite line l and any point P not on l, there does not exist a line which passes through P and is parallel to l. Hyperbolic geometry is also known as saddle geometry or Lobachevskian geometry. It differs in many ways to Euclidean geometry, often leading to quite counter-intuitive results. Some of these remarkable consequences of this geometry's unique fifth postulate include: 1. The sum of the three interior angles in a triangle is strictly less than 180°. Moreover, the angle sums of two distinct triangles are not necessarily the same. 2. Two triangles with the same interior angles have the same area. Models of Hyperbolic Space The following are four of the most common models used to describe hyperbolic space. 1. The Poincaré Disc Model. Also known as the conformal disc model. In it, the hyperbolic plane is represented by the interior of a circle, and lines are represented by arcs of circles that are orthogonal to the boundary circle and by diameters of the boundary circle. Preserves hyperbolic angles. 2. The Klein Model. Also known as the Beltrami-Klein model or projective disc model. In it, the hyperbolic plane is represented by the interior of a circle, and lines are represented by chords of the circle. This model gives a misleading visual representation of the magnitude of angles. 3. The Poincaré Half-Plane Model. The hyperbolic plane is represented by one-half of the Euclidean plane, as defined by a given Euclidean line l, where l is not considered part of the hyperbolic space. Lines are represented by half-circles orthogonal to l or rays perpendicular to l. Preserves hyperbolic angles. 4. The Lorentz Model. Spheres in Lorentzian four-space. The hyperbolic plane is represented by a two-dimensional hyperboloid of revolution embedded in three-dimensional Minkowski space. Based on this geometry's definition of the fifth axiom, what does parallel mean? The following definitions are made for this geometry. If a line l and a line m do not intersect in the hyperbolic plane, but intersect at the plane's boundary of infinity, then l and m are said to be parallel. If a line p and a line q neither intersect in the hyperbolic plane nor at the boundary at infinity, then p and q are said to be ultraparallel. The Ultraparallel Theorem For any two lines m and n in the hyperbolic plane such that m and n are ultraparallel, there exists a unique line l that is perpendicular to both m and n. Elliptic geometry differs in many ways to Euclidean geometry, often leading to quite counter-intuitive results. For example, directly from this geometry's fifth axiom we have that there exist no parallel lines. Some of the other remarkable consequences of the parallel postulate include: The sum of the three interior angles in a triangle is strictly greater than 180°. Models of Elliptic Space Spherical geometry gives us perhaps the simplest model of elliptic geometry. Points are represented by points on the sphere. Lines are represented by circles through the points. - Euclid's First Four Postulates - Euclid's Fifth Postulate - Incidence Geometry - Projective and Affine Planes (necessary?) - Axioms of Betweenness - Pasch and Crossbar - Axioms of Congruence - Continuity (necessary?) - Hilbert Planes - Neutral Geometry If you would like to request anything in this topic please post it below. - Modern geometry - An Alternative Way and Alternative Geometric Means of Calculating the Area of a Circle = Geometry/An Alternative Way and Alternative Geometric Means of Calculating the Area of a Circle

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Tsunamis can occur anywhere along the United States coastline, but are most likely to strike along the coasts of Alaska, California, Oregon, Washington, and Hawaii. Tsunamis, sometimes referred to as seismic sea waves or incorrectly as tidal waves, are abnormally large ocean waves that can reach heights of 100 feet tall and travel at speeds of up to 500 miles per hour. They are usually triggered by underground volcanic activity and earthquakes, but can also occur from undersea landslides or even nuclear explosions. Tsunamis are difficult to detect at sea because they are only a few inches high. As a tsunami approaches land, the height of the wave increases while at the same time, the water at the shoreline recedes. While tsunamis are difficult to detect, there are systems in place that provide some early indications that one is eminent or on its way. When a tsunami is detected, or if a possible trigger even occurs, a warning is issued to areas impacted by a possible strike. Advisory - An earthquake occurred in the Pacific Ocean, which might trigger a tsunami. Watch - A tsunami was triggered, but is two hours or more travel time away from the area in Watch status. Warning - A tsunami is most likely on its way that can cause damage; People in the warned area are advised evacuate. Being prepared for a tsunami or any other type of disaster requires preparation beforehand. There are three important steps you should complete in preparation for a tsunami. To make this even easier, Essential Packs provides you with a FREE online Emergency Preparedness Planning Guide that takes your step-by-step through the process: Having the proper emergency preparedness kit, having a plan, and knowing what to do before a tsunami strikes, will greatly improve you and your family’s chances of being safe. Complete the 3-Step Emergency Preparedness Planning Guide today! Getting an Emergency Preparedness Kit for your home, office (or school), and car(s) is an essential first step in being prepared for a tsunami. Emergency Preparedness Kits from Essential Packs, provide you and your family with the emergency supplies your family needs to last for 3 days (72 hours). Deluxe Kits from Essential Packs are compliant with FEMA's guidelines and include important items like: emergency food and drinking water, Flashlight, radio, first-aid supplies, sanitation supplies, emergency blankets, waterproof ponchos, and much more. For additional help on selecting the right kit, visit Step 1 - Get A Kit of our Emergency Preparedness Planning Guide. In order to know what to do when a tsunami occurs, you need to create a Family Emergency Plan. Sit down with your family members and decide how you will get in contact with each other, where you will go, and what you will do in the event of a disaster or emergency. To make this easy, Essential Packs provides you with a Family Emergency Planning Document that you can download for FREE. Simply open this PDF document and fill-in the blanks, then, print a copy for each family member, and store one copy in your Emergency Preparedness Kit. You should update your Family Emergency Plan every six months, as phone number, work locations, and other important information could change. For more help on creating a Family Emergnecy Plan, visit Step 2 - Make A Plan of our Emergency Preparedness Planning Guide. The final step to getting prepared is to be informed about what to do before, during, and after a tsunami. FEMA's In-Depth Citizen's Guide to Disaster Preparedness helps you do this by providing you with comprehensive emergency preparedness information a variety of disasters. Visit Step 3 - Be Informed of our Emergency Preparedness Planning Guide to download FEMA's comprehensive, 200 page book called, "Are You Ready? An In-Depth Guide to Citizen Preparedness".

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Equal Rights Amendment |Part of the series on| |U.S. Discrimination Law| |Standards of Review| |Other Legal Theories| |Defining Moments in Law| |Modalities of Constitutional Law| The so-called Equal Rights Amendment (ERA) was a proposed amendment to the United States Constitution passed by Congress in 1972 and sent to the states for ratification with a deadline for ratification of seven years. It was stopped by a conservative grass-roots movement that raised a number of objections, including the following: - it would require drafting women just like men, and putting women in combat just like men - it would require taxpayer-funded abortion - it would require same-sex marriage - it would expand the power of the federal government and the courts - it would require identical treatment of men and women, and boys and girls, in many other areas of life, such as schools and school activities The amendment, which was untitled, stated: - SECTION 1. Equality of rights under the law shall not be denied or abridged by the United States or by any State on account of sex. - SECTION 2. The Congress shall have the power to enforce, by appropriate legislation, the provisions of this article. - SECTION 3. This amendment shall take effect two years after the date of ratification. The purpose of the ERA was to prohibit many legal distinctions between men and women, and boys and girls. An all-male draft or male-only combat regiments would presumably have become illegal. But equal representation of boys and girls on sports teams, and the termination of all-boys or all-girls sports teams, would probably not have been required, since the courts usually treat some gender distinctions as a bona fida operating qualification for which discrimination is allowed. Thirty-five out of the required 38 states ratified the amendment, but opposition led by Phyllis Schlafly ultimately defeated it. Congress then extended the deadline to 1982 in legislation that was later invalidated, though no other states ratified the amendment in the additional three years. Several states rescinded their prior ratification of the amendment. After women won the right to vote in 1920, the next step for some feminists was the the ERA; it was introduced in Congress and rejected in every one of its sessions from 1923 until 1972. It was initially authored by Alice Paul, head of the National Women's Party, who led the suffrage campaign. The original version of the ERA stated that "Men and women shall have equal rights throughout the United States and every place subject to its jurisdiction." The ERA was opposed by women who wanted special laws to protect women workers, including Eleanor Roosevelt. It was mostly supported by middle class Republican women. The ERA was advocated for by Gloria Steinem and Betty Friedan. For example, a pro-ERA site claims that "[t]he new Constitution's promised rights were fully enjoyed only by certain white males. Women were treated according to social tradition and English common law and were denied most legal rights. In general they could not vote, own property, keep their own wages, or even have custody of their children." In fact, the U.S. Constitution actively denied women the right to vote in the late 1800s by including gender specific language in the Fourteenth Amendment. Section two of the Fourteenth Amendment specifically uses the phrase "male inhabitants" to continue the disenfranchisment of women. The Canadian equivalent to the ERA is known as Charter Section 28. It was ratified as a section of the Canadian Constitution in 1982, and came into effect in 1985. Since then, many laws have been invalidated because they were in violation of Charter Section 28.

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Delegates from twelve states met in Philadelphia in 1787 to discuss strengthening the powers of the national government. Only Rhode Island failed to send delegates to the convention. After more than four months of secret debates, proposals, and compromises, a majority of the delegates reached an agreement. The Constitution was signed by thirty-nine of the fifty-five members of the Constitutional Convention on September 17, 1787, in Philadelphia. The delegation from Virginia played an integral role throughout the Constitutional Convention. Virginia was the first state to agree to send delegates to Philadelphia. Governor Edmund Randolph introduced a plan, authored by James Madison, to dispose of the Articles of Confederation and to establish a fundamentally different form of government. The Virginia Plan served as the foundation for the newly written Constitution. George Washington served as president of the Convention, and lent his reputation to the proceedings. Madison's tireless work and leadership behind the scenes earned him the nickname “The Father of the Constitution.” Another Virginia delegate, John Blair, was a supporter of the Constitution, but had more influence after the Convention, as a member of the first Supreme Court of the United States. George Mason contributed greatly to the debate and formation of the Constitution, but along with Randolph ultimately refused to sign in part because the Constitution lacked a Declaration of Rights. George Wythe was a signer of the Declaration of Independence, and the first law professor in America. Wythe was assigned, along with Alexander Hamilton, to create the rules and procedures that governed the Constitutional Convention. James McClurg was added as a delegate after Patrick Henry, Thomas Nelson, and Richard Henry Lee refused the appointment. Like Mason and Randolph, neither McClurg nor Wythe signed the final Constitution; Wythe had to leave the Convention early and eventually voted for ratification at the Virginia Convention. Therefore, despite significant contributions, only three of the seven Virginians in Convention signed the final copy of the Constitution. Congress submitted the draft Constitution to the states for ratification. In a letter to Congress that accompanied the Constitution, Washington, who served as president of the Convention, addressed the difficulties of the Convention, with “a difference among the several States as to their situation, extent, habits, and particular interests … we kept steadily in our view that which appears to us the greatest interest of every true American, the consolidation of our Union, in which is involved our prosperity, felicity, safety—perhaps our national existence.” This copy of the Constitution may have been the first one made for public consumption. It is believed that this printing took place only days after the Constitutional Convention members reached their final agreement. Copies were distributed all throughout the states, as citizens and delegates to the ratification conventions looked to learn as much about the proposed system of government as possible. On June 21, 1788, New Hampshire became the ninth state to ratify the Constitution, which made the adoption official. The conventions in Virginia and New York ratified soon after. On July 2, 1788, the Continental Congress announced the adoption of the Constitution, and on October 10, the Continental Congress officially ended its business. 1. How many articles are there in the Constitution? 2. Who was the president of the Constitutional Convention? 3. What were some of the issues and differences that existed among the delegates of the Constitutional Convention? How were they resolved? 4. The final copy of the Constitution was shorter than all of the drafts. Why do you think the signers made that decision? What benefits might a shorter, more vague document present over a longer more detailed one? What disadvantages? 1. Compare this finalized copy of the Constitution with the Virginia Plan. Where were there changes made? Why do you think these changes were made? Do you think that the changes are significant? How do the changes affect your understanding and the way you think about the Constitution? 2. When compared with the Articles of Confederation what features of the U.S. Constitution represent an improvement over the previous form of government? 3. The United States Constitution is often called a “living document.” What does this mean to you? Do you agree with that classification being applied to the Constitution? This printing of the Constitution is described in Prologue, the Journal of the National Archives (Fall, 1970): 82. There it is suggested that this printing followed immediately the printing in Dunlap and Claypoole's Pennsylvania Packet on Wednesday, Sept. 19, 1787. This Library of Congress copy of this broadside is damaged with some loss of text. Kaminski, John P., and Gaspare J. Saladino, eds. The Documentary History of the Ratification of the Constitution, Vol 8: Ratification of the Constitution by the States: Virginia, Vol. I:xxxv–xxxvi. Madison: State Historical Society of Wisconsin, 1988. Briceland, Alan V. “Virginia: The Cement of the Union.” In The Constitution and the States: The Role of the Original Thirteen in the Framing and Adoption of the Federal Constitution. Edited by Patrick T. Conley and John P. Kaminski. Madison, Wis.: Madison House, 1988. Gales, Joseph. The Debates and Proceedings in the Congress of the United States (Annals of Congress). Vol. 1, Introduction, v–xiii. Washington, D.C.: Gales and Seaton, 1834.

http://www.virginiamemory.com/online_classroom/shaping_the_constitution/doc/constitution

Learn something new every day More Info... by email On 16 December 1773, colonists in Boston, Massachusetts, boarded three ships full of tea sent by the East India Company and threw the tea into the Boston Harbor. They did so to protest the Tea Act, passed by the Parliament of Great Britain, which practically granted the East India Company a monopoly on tea distribution in the American colonies. The colonists were required to pay a tax on this tea because of the Townshend Acts. One of the first effects of the Boston Tea Party was the passage of strict new laws known as the "Coercive Acts," which were intended to punish the colonists. This led to many of the later effects, including the formation of the First Continental Congress. The Boston Tea Party was also a major catalyst leading to the start of the American Revolutionary War. The Coercive Acts, also known as the Intolerable Acts, were intended as punishment for the destruction of property. One of these acts actually shut down trade in the city of Boston until the damages could be repaid. These punitive acts were particularly strict because the British Parliament intended to make an example of Boston in order to discourage other acts of rebellion from the colonies. These acts led to another of the effects of the Boston Tea Party: the development of the First Continental Congress. This congress was formed with the intent of petitioning the British Parliament in order to have the harsh punitive acts repealed. The petition was ignored, leading to the development of the Second Continental Congress. Between these two Congresses, the American Revolutionary War began, and the Second Continental Congress was highly involved in overseeing the war effort. In many ways, the American Revolutionary War itself can be seen as one of the effects of the Boston Tea Party. There were many different factors and rebellious acts that eventually precipitated in the war, but the Boston Tea Party was among the most significant. The passage of the Coercive Acts and the subsequent refusal of the British Parliament to repeal them led to a great deal of dissatisfaction from colonists. It is distinctly possible that without the Boston Tea Party, the Revolution would have been delayed for quite some time or may not have even happened at all. This is a great source and it really helped me have an easier time with my research! I'm glad I found WiseGeek!

http://www.wisegeek.org/what-were-the-effects-of-the-boston-tea-party.htm

Being able to read all the information presented on the Periodic Table is the most important chemistry skill you will learn in this class. A Periodic Table will always be available to you, in this class and even in college classes. If you can use that information, you have almost everything you need right in front of you. Paper Periodic Table Assignment: - Print this table for class notes. - Add atomic mass, rounded to the nearest whole number. - Add element names. Make notes on the front and back and use this "Paper Periodic Table" on any class test. If it is lost, you may not make a photocopy of another student's table. You may print another copy and re-write your notes onto it. Students must have their OWN Paper Periodic Table, or use the one on the classroom wall. No other notes or computers may be used on "closed-computer" tests. A link to the class online Periodic Table is found on the lower right of any chemistry class concept page. Roll over a square to see the element name. Click on a square for detailed information about the element.

http://crescentok.com/staff/jaskew/ISR/TigerChem/periodic_table/table2.htm

Johannes Kepler (1571-1630), born in Weil der Stadt, Württemburg, the Holy Roman Empire of German Nationality. Kepler eventually moved to Prague during the Thirty-Years War to escape religious persecution. While in Prague, he worked with Tycho Brahe, a renowned Danish astronomer. After Brahe passed away, Kepler followed in his footsteps; using Brahe’s collection of data, Kepler was the first to correctly explain planetary motion. Kepler’s three laws of planetary motion that bear his name were published in 1609 and 1619. Kepler described how the orbits of the planets were not the circles described by Aristotle and assumed implicitly by Copernicus, but were instead "flattened circles", or ellipses. Kepler also made important contributions to optics, including the first explanation of the human vision process by refraction within the eye, the first eyeglass design for nearsightedness and farsightedness, the first description of depth perception, and the first detailed explanation of the principles of how a telescope works. Kepler also theorized that the tides were caused by the Earth's moon, and that the Sun rotates about its axis. Johannes Kepler: His Life, His Laws and Times (NASA's Kepler Mission Website) Kepler Biography (University of St. Andrews, Scotland, School of Mathematics and Statistics) NASA's Kepler Mission: a search for habitable planets

http://www.eoearth.org/article/Kepler,_Johannes

Home » Education » Special Education Straightforward, practical, and user friendly, this unique guide addresses an essential component of decision making in schools. The authors show how systematic screenings of behavior—used in conjunction with academic data—can enhance teachers' ability to teach and support all students within a response-to-intervention framework. Chapters review reliable, valid screening measures for all grade levels, discuss their strengths and weaknesses, and explain how to administer, score, and interpret them. Practitioners get helpful guidance for evaluating their school's needs and resources and making sound choices about which tools to adopt.

http://www.guilford.com/cgi-bin/cartscript.cgi?page=pr/lane3.htm&dir=edu/speced&cart_id=877143.5617

Parts of the cell The cell membrane is very important, because it works as a selective filter that allows only certain things to come inside or go outside the cell, it act as a body guard for our body.It can maintain a stable and healthy environment for cell in order to keep people healthy. plant cell membranes are rigid walls, and animal cell membranes are lipid bilayers. The phospholipid bilayer which the cell membrane is an example of, is composed of various cholesterol, phospholipids, glycolipids, blagoscony and proteins. Below is an example of a simple phospholipid bilayer. The smaller molecules shown between the phospholipids are Cholesterol molecules. They help to provide rigidity or stability to the membrane. The two main components of phospholipids are shown in these figures by blue circles representing the hydrophilic head groups and by long thin lines representing the hydrophobic fatty acid tails. Both the interior of the cell and the area surrounding the cell is made up of water or similar aqueous solution. Consequently, phospholipids orient themselves with respect to the water and with each other so that the hydrophilic ("water loving") head groups are grouped together and face the water, and the hydrophobic ("water fearing") tails turn away from the water and toward each other. This self-organization of phospholipids results in one of just a few easily recognizable structures. Cell membranes are constructed of a phospholipid bilayer as shown above. Smaller structures can also form, known as 'micelles' in which there is no inner layer of phospholipid. Instead, the interior of a micell is wholly hydrophobic, filled with the fatty acid chains of the phospholipids and any other hydrophobic molecule they enclose. Micelles are not so important for the understanding of cellular structure, but are useful for demonstrating the principles of hydrophilicity and hydrophobicity, and for contrasting with lipid bilayers. At least 10 different types of lipids are commonly found in cell membranes. Each type of cell or organelle will have a different percentage of each lipid, protein and carbohydrate. The main types of lipids are: - Diphosphatidylglycerol (Cardiolipin) - Phosphatidic acid The Cell Membrane is Asymmetric The cell membrane tends to have different composition on one side of the membrane than on the other side of the membrane. The differences can be caused by the different ratios or types of amphipathic lipid-based molecules, the different positioning of the proteins (facing in or facing out), or the fixed orientations of proteins spanning the membrane. Additionally, there are different enzymatic activities in the outer and inner membrane surfaces. The reason the cell membrane is asymmetric is because when the proteins are synthesized by the preexisting membranes, they are inserted into the membrane in an asymmetric manner. The asymmetry of the cell membrane allows the membrane to be rigid and allows the cell to have a different intracellular environment from the existing extracellular environment. Additionally, the cell membrane's phospholipids are distributed asymmetrically across the lipid bilayer, in a phenomenon called membrane phospholipid asymmetry. There are three mechanisms for transmembrane movement of phospholipids: 1) spontaneous diffusion, 2) facilitated diffusion, 3) ATP-dependent active translocation. The spontaneous diffusion is a form of passive transport. Because passive transport does not require energy to transport non-polar substances through the membrane, this can happen spontaneously. Facilitated diffusion, like spontaneous diffusion, is a form of passive transport. The molecules or ions in this diffusion pass through the membrane by using specific transmembrane transport proteins. Membrane transport of small molecules Because animal membrane proteins are lipid bilayer which are inner hydrophobic, this character prohibits polar molecules. Transport proteins can provide help for this situation. It can transport polar molecules across the membrane. There are several types of membrane transport proteins. They are uniports and cotransport. Uniports can move solutes from one side to another, change the position of the proteins. Cotransport systems can simultaneously sending two solutes across the lipid bilayer. Solutes are sent in the same direction or opposite directions Transport proteins does not need to be acts natural direction. Membrane Transport of Macromolecules Membrane transport of Macromolecules can divide into two parts, they are exocytosis and endocytosis. In exocytosis, the contents of vesicles are released when the vesicle fuses with the cell membrane. There are five steps involved, which are vesicle trafficking, vesicke tethering, vesicle docking, vesicle priming and vesicle fusion. In endocytosis the membrane depresses and pinches off, enclosing the molecule. In receptor-mediated endocytosis, coated pits and vesicles bind to specific receptors on the cell surface, allowing the cell to select what molecules to take and what to reject.

http://en.m.wikibooks.org/wiki/Cell_Biology/Membranes

A theorem generally has a set-up - a number of conditions, which may be listed in the theorem or described beforehand. Then it has a conclusion - a mathematical statement which is true under the given set up. The proof, though necessary to the statement's classification as a theorem is not considered part of the theorem. In general mathematics a statement must be interesting or important in some way to be called a theorem. Less important statements are called: - lemma: a statement that forms part of the proof of a larger theorem. Of course, the distinction between theorems and lemmas is rather arbitrary, since one mathematician's major result is another's minor claim. Gauss' Lemma and Zorn's Lemma, for example, are interesting enough per se for some authors to stop at the nominal lemma without going on to use that result in any "major" theorem. - corollary: a statement which follows immediately or very simply from a theorem. A proposition A is a corollary of a proposition or theorem B if A can be deduced quickly and easily from B. - proposition: a result not associated with any particular theorem. - claim: a very minor, but necessary or interesting result, which may be part of the proof of another statement. Despite the name, claims are proven. - remark: similar to claim. Probably presented without proof, which is assumed to be obvious. As noted above, a theorem requires some sort of logical framework, this will consist of a basic set of axioms (see axiomatic system), as well as a process of inference, which allows to derive new theorems from axioms and other theorems that have been derived earlier. In propositional logic, any proven statement is called a theorem. - mathematics for a list of famous theorems and conjectures. - Gödel's incompleteness theorem

http://www.encyclopedia4u.com/t/theorem.html

The Origin of Asteroidsby Dr. Walt Brown (This article has been reproduced with permission from the Center for Scientific Creation. The original article can be found here.) NOTE - In order to fully understand the content of this article (and it’s companion article The Origin of Comets), you should read the book, In the Beginning by Dr. Walt Brown. This book fully explains Dr. Brown’s Hydroplate Theory which is the foundation upon which this article is written. In fact, this “article” is actually a chapter in the book, In the Beginning. Members of the 4th Day Alliance can download the complete PDF copy of this chapter by clicking here. Figure 156: Asteroid Ida and Its Moon, Dactyl. In 1993, the Galileo spacecraft, heading toward Jupiter, took this picture 2,000 miles from asteroid Ida. To the surprise of most, Ida had a moon (about 1 mile in diameter) orbiting 60 miles away! Both Ida and Dactyl are composed of earthlike rock. We now know of 68 other asteroids that have moons.1 According to the laws of orbital mechanics (described in the preceding chapter), capturing a moon in space is unbelievably difficult—unless both the asteroid and a nearby potential moon had very similar speeds and directions and unless gases surrounded the asteroid during capture. If so, the asteroid, its moon, and each gas molecule were probably coming from the same place and were launched at about the same time. Within a million years, passing bodies would have stripped the moons away, so these asteroid-moon captures must have been recent. From a distance, large asteroids look like big rocks. However, many show, by their low density, that they contain either much empty space or something light, such as water ice.2 Also, the best close-up pictures of an asteroid show millions of smaller rocks on its surface. Therefore, asteroids are flying rock piles held together by gravity. Ida, about 35 miles long, does not have enough gravity to squeeze itself into a spherical shape. SUMMARY: The fountains of the great deep launched rocks as well as muddy water. As rocks moved farther from Earth, Earth’s gravity became less significant to them, and the gravity of nearby rocks became increasingly significant. Consequently, many rocks, assisted by their mutual gravity and surrounding clouds of water vapor, merged to become asteroids. Isolated rocks in space are meteoroids. Drag forces caused by water vapor and thrust forces produced by the radiometer effect concentrated asteroids in what is now the asteroid belt. All the so-called “mavericks of the solar system” (asteroids, meteoroids, and comets) resulted from the explosive events at the beginning of the flood. Asteroids, also called minor planets, are rocky bodies orbiting the Sun. The orbits of most asteroids lie between those of Mars and Jupiter, a region called the asteroid belt. The largest asteroid, Ceres, is almost 600 miles in diameter and has about one-third the volume of all other asteroids combined. Orbits of almost 30,000 asteroids have been calculated. Many more asteroids have been detected, some less than 20 feet in diameter. A few that cross the Earth’s orbit would do great damage if they ever collided with Earth. Two explanations are given for the origin of asteroids: (1) they were produced by an exploded planet, and (2) a planet failed to evolve completely. Experts recognize the problems with each explanation and are puzzled. The hydroplate theory offers a simple and complete—but quite different—solution that also answers other questions. Meteorites, Meteors, and MeteoroidsIn space, solid bodies smaller than an asteroid but larger than a molecule are called “meteoroids.” They are renamed “meteors” as they travel through Earth’s atmosphere, and “meteorites” if they hit the ground. Exploded-Planet Explanation. Smaller asteroids are more numerous than larger asteroids, a pattern typical of fragmented bodies. Seeing this pattern led to the early belief that asteroids are remains of an exploded planet. Later, scientists realized that all the fragments combined would not make up one small planet.3 Besides, too much energy is needed to explode and scatter even the smallest planet. Failed-Planet Explanation. The most popular explanation today for asteroids is that they are bodies that did not merge to become a planet. Never explained is how, in nearly empty space, matter merged to become these rocky bodies in the first place,4 why rocky bodies started to form a planet but stopped,5 or why it happened only between the orbits of Mars and Jupiter. Also, because only vague explanations have been given for how planets formed, any claim to understand how one planet failed to form lacks credibility. In general, orbiting rocks do not merge to become either planets or asteroids. Special conditions are required, as explained on page 267 and Endnote 23 on page 288.] Today, collisions and near collisions fragment and scatter asteroids, just the opposite of this “failed-planet explanation.” In fact, during the 4,600,000,000 years evolutionists say asteroids have existed, asteroids would have had so many collisions that they should be much more fragmented than they are today.6 Hydroplate Explanation. Asteroids are composed of rocks expelled from Earth. The size distribution of asteroids does show that at least part of a planet fragmented. Although an energy source is not available to explode and disperse an entire Earth-size planet, the eruption of so much supercritical water from the subterranean chambers could have launched one 2,300th of the Earth—the mass of all asteroids combined. Astronomers have tried to describe the exploded planet, not realizing they were standing on the remaining 99.95% of it—too close to see it.7 As flood waters escaped from the subterranean chambers, pillars, forced to carry more and more of the weight of the overlying crust, were crushed. Also, the almost 10-mile-high walls of the rupture were unstable, because rock is not strong enough to support a cliff more than 5 miles high. As lower portions of the walls were crushed, large blocks8 were swept up and launched by the jetting fountains. Unsupported rock in the top 5 miles then fragmented. The smaller the rock, the faster it accelerated and the farther it went, just as a rapidly flowing stream carries smaller dirt particles faster and farther. Water droplets in the fountains partially evaporated and quickly froze. Large rocks had large spheres of influence which grew as the rocks traveled away from Earth. Larger rocks became “seeds” around which other rocks and ice collected as spheres of influence expanded. Because of all the evaporated water vapor and the resulting aerobraking, even more mass concentrated around the “seeds.”Clumps of rocks became asteroids. Question 1: Why did some clumps of rocks and ice in space become asteroids and others become comets?Imagine living in a part of the world where heavy frost settled each night, but the Sun shone daily. After many decades, would the countryside be buried in hundreds of feet of frost? The answer depends on several things besides the obvious need for a large source of water. If dark rocks initially covered the ground, the Sun would heat them during the day, so frost from the previous night would tend to evaporate. However, if the sunlight was dim or the frost was thick (thereby reflecting more sunlight during the day), little frost would evaporate. More frost would accumulate the next night. Frost thickness would increase every 24 hours. Now imagine living on a newly formed asteroid. Its spin would give you day-night cycles. After sunset, surface temperatures would plummet toward nearly absolute zero (-460°F), because asteroids do not have enough gravity to hold an atmosphere for long. With little atmosphere to insulate the asteroid, the day’s heat would quickly radiate, unimpeded, into outer space. Conversely, when the Sun rose, its rays would have little atmosphere to warm, so temperatures at the asteroid’s surface would rise rapidly. As the fountains of the great deep launched rocks and water droplets, evaporation in space dispersed an “ocean” of water molecules and other gases in the inner solar system. Gas molecules that struck the cold side of your spinning asteroid would become frost.9 Sunlight would usually be dim on rocks in larger, more elongated orbits. Therefore, little frost would evaporate during the day, and the frost’s thickness would increase. Your “world” would become a comet. However, if your “world” orbited relatively near the Sun, its rays would evaporate each night’s frost, so your “world” would remain an asteroid. Heavier rocks could not be launched with as much velocity as smaller particles (dirt, water droplets, and smaller rocks). The heavier rocks merged to become asteroids, while the smaller particles, primarily water, merged to become comets, which generally have larger orbits. No “sharp line” separates asteroids and comets. PREDICTION 33Asteroids are rock piles, often with ice acting as a weak “glue” inside. Large rocks that began the capture process are nearer the centers of asteroids. Comets, which are primarily ice, have rocks in their cores. Four years after this prediction was published in 2001 (In the Beginning, 7th edition, page 220), measurements of the largest asteroid, Ceres, found that it does indeed have a dense, rocky core and primarily a water-ice mantle.10 Question 2: Wasn’t asteroid Eros found to be primarily a large, solid rock?A pile of dry sand here on Earth cannot maintain a slope greater than about 30 degrees. If it were steeper, the sand grains would roll downhill. Likewise, a pile of dry pebbles or rocks on an asteroid cannot have a slope exceeding about 30 degrees. However, 4% of Eros’ surface exceeds this slope, so some scientists concluded that much of Eros must be a large, solid rock. This conclusion overlooks the possibility that ice is present between some rocks and acts as a weak glue—as predicted above. Ice in asteroids would also explain their low density. Endnote 8 gives another reason why asteroids are probably flying rock piles. Question 3: Objects launched from Earth should travel in elliptical, cometlike orbits. How could rocky bodies launched from Earth become concentrated in almost circular orbits between Mars and Jupiter?Gases, such as water vapor and its components,11 were abundant in the inner solar system for many years after the flood. Hot gas molecules striking each asteroid’s hot side were repelled with great force. This jetting action was like air rapidly escaping from a balloon, applying a thrust in a direction opposite to the escaping gas.12 Cold molecules striking each asteroid’s cold side produced less jetting. This thrusting, efficiently powered by solar energy, pushed asteroids outward, away from the sun, concentrating them between the orbits of Mars and Jupiter.13 [See Figures 157 and 158.] Figure 157: Thrust and Drag Acted on Asteroids (Sun, asteroid, gas molecules, and orbit are not to scale.) The fountains of the great deep launched rocks and muddy water from Earth. The larger rocks, assisted by water vapor and other gases within the spheres of influence of these rocks, captured other rocks and ice particles. Those growing bodies that were primarily rocks became asteroids. The Sun heats an asteroid’s near side, while the far side radiates its heat into cold outer space. Therefore, large temperature differences exist on opposite sides of each rocky, orbiting body. The slower the body spins, the darker the body,14 and the closer it is to the Sun, the greater the temperature difference. (For example, temperatures on the sunny side of our Moon reach a searing 260°F, while on the dark side, temperatures can drop to a frigid -280°F.) Also, gas molecules (small blue circles) between the Sun and asteroid, especially those coming from very near the Sun, are hotter and faster than those on the far side of an asteroid. Hot gas molecules hitting the hot side of an asteroid bounce off with much higher velocity and momentum than cold gas molecules bouncing off the cold side. Those impacts slowly expanded asteroid orbits until too little gas remained in the inner solar system to provide much thrust. The closer an asteroid was to the Sun, the greater the outward thrust. Gas molecules, densely concentrated near Earth’s orbit, created a drag on asteroids. My computer simulations have shown how gas, throughout the inner solar system for years after the flood, herded asteroids into a tight region near Earth’s orbital plane—an asteroid belt.15 Thrust primarily expanded the orbits. Drag circularized orbits and reduced their angles of inclination. Figure 158: The Radiometer Effect. This well-known novelty, called a radiometer, demonstrates the unusual thrust that pushed asteroids into their present orbits. Sunlight warms the dark side of each vane more than the light side. The partial vacuum inside the bulb approaches that found in outer space, so gas molecules travel relatively long distances before striking other molecules. Gas molecules bounce off the hotter, black side with greater velocity than off the colder, white side. This turns the vanes away from the dark side. The black side also radiates heat faster when it is warmer than its surroundings. This can be demonstrated by briefly placing the radiometer in a freezer. There the black side cools faster, making the white side warmer than the black, so the vanes turn away from the white side. In summary, the black side gains heat faster when in a hot environment and loses heat faster when in a cold environment. Higher gas pressure always pushes on the warmer side. Question 4: Could the radiometer effect push asteroids 1–2 astronomical units (AU) farther from the Sun?Each asteroid began as a swarm of particles (rocks, ice, and gas molecules) orbiting within a large sphere of influence. Because a swarm’s volume was quite large, its spin was much slower than it would be as it shrank to become an asteroid—perhaps orders of magnitude slower. The slow spin produced extreme temperature differences between the hot and cold sides. The cold side would have been so cold that gas molecules striking it would tend to stick, thereby adding “fuel” to the developing asteroid. Because the swarm’s volume was large, the radiometer pressure acted over a large area and produced a large thrust. The swarm’s large thrust and low density caused the swarm to rapidly accelerate—much like a feather placed in a gentle breeze. Also, the Sun’s gravity 93,000,000 miles from the Sun (the Earth-Sun distance) is 1,600 times weaker than Earth’s gravity here on Earth.17 So, pushing a swarm of rocks and debris farther from the Sun was surprisingly easy, because there is almost no resistance in outer space. Question 5: Why are 4% of meteorites almost entirely iron and nickel? Also, why do meteorites rarely contain quartz, which constitutes about 27% of granite’s volume?Pillars were formed in the subterranean chamber when the thicker portions of the crust were squeezed downward onto the chamber floor. Twice daily, during the centuries before the flood, these pillars were stretched and compressed by tides in the subterranean water. This gigantic heating process steadily raised pillar temperatures. [See “What Triggered the Flood?” here.] As explained in Figure 159, temperatures in what are now iron-nickel meteorites once exceeded 1,300°F, enough to dissolve quartz and allow iron and nickel to settle downward and become concentrated in the pillar tips.18 (A similar gravitational settling process concentrated iron and nickel in the Earth’s core after the flood began. See “Melting the Inner Earth” here.) Evolutionists have great difficulty explaining iron-nickel meteorites. First, everyone recognizes that a powerful heating mechanism must first melt at least some of the parent body from which the iron-nickel meteorites came, so iron and nickel can sink and be concentrated. How this could have occurred in the weak gravity of extremely cold asteroids has defied explanation.19 Second, the concentrated iron and nickel, which evolutionists visualize in the core of a large asteroid, must then be excavated and blasted into space. Available evidence shows that this has not happened.20 Figure 156: Asteroid Ida and Its Moon, Dactyl. Most iron-nickel meteorites display Widmanstätten patterns. That is, if an iron-nickel meteorite is cut and its face is polished and then etched with acid, the surface has the strange crisscross pattern shown above. This shows that temperatures throughout those meteorites exceeded 1,300°F.16 Why were so many meteoroids, drifting in cold space, at one time so uniformly hot? An impact would not produce such uniformity, nor would a blowtorch. The heating a meteor experiences in passing through the atmosphere is barely felt more than a fraction of an inch beneath the surface. If radioactive decay generated the heat, certain daughter products should be present; they are not. Question 5 explains how these high temperatures were probably reached. Question 6: Aren’t meteoroids chips from asteroids?This commonly-taught idea is based on an error in logic. Asteroids and meteoroids have some similarities, but that does not mean that one came from the other. Maybe a common event produced both asteroids and meteoroids. Also, three major discoveries suggest that meteoroids came not from asteroids, but from Earth. 1. In the mid-1970s, the Pioneer 10 and 11 spacecraft traveled out through the asteroid belt. NASA expected that the particle detection experiments on board would find 10 times more meteoroids in the belt than are present near Earth’s orbit.21 Surprisingly, the number of meteoroids diminished as the asteroid belt was approached.22 This showed that meteoroids are not coming from asteroids but from nearer the Earth’s orbit. 2. A faint glow of light, called the zodiacal light, extends from the orbit of Venus out to the asteroid belt. The light is reflected sunlight bouncing off dust-size particles. This lens-shaped swarm of particles orbits the Sun, near Earth’s orbital plane. (On dark, moonless nights, zodiacal light can be seen in the spring in the western sky after sunset and in the fall in the eastern sky before sunrise.) Debris chipped off asteroids would have a wide range of sizes and would not be as uniform and fine as the particles reflecting the zodiacal light. Debris expelled by the fountains of the great deep would place fine dust particles in the Earth's orbital plane. 3. Many meteorites have remanent magnetism, so they must have come from a larger magnetized body. Eros, the only asteroid on which a spacecraft has landed and taken magnetic measurements, has no net magnetic field. If this is true of other asteroids as well, meteorites probably did not come from asteroids.30 If asteroids are flying rock piles, as it now appears, any magnetic fields in the randomly oriented rocks would be largely self-canceling, so the asteroid would have no net magnetic field. Therefore, instead of coming from asteroids, meteorites likely came from a magnetized body such as a planet. Because Earth’s magnetic field is 2,000 times greater than that of all other rocky planets combined, meteorites probably came from Earth. Remanent magnetism decays, so meteorites must have recently broken away from their parent magnetized body. Those who believe that meteorites were chipped off asteroids say this happened millions of years ago. PREDICTION 34:Most rocks comprising asteroids will be found to be magnetized. Two InterpretationsWith a transmission electron microscope, Japanese scientist Kazushige Tomeoka identified several major events in the life of one meteorite. Initially, this meteorite was part of a much larger parent body orbiting the Sun. The parent body had many thin cracks, through which mineral-rich water cycled. Extremely thin mineral layers were deposited on the walls of these cracks. These deposits, sometimes hundreds of layers thick, contained calcium, magnesium, carbonates, and other chemicals. Mild thermal metamorphism in this rock shows that temperatures increased before it experienced some final cracks and was blasted into space.31 Hydroplate Interpretation. Earth was the parent body of all meteorites, most of which came from pillars. [Pages 381–386 explain how, why, when, and where pillars formed.] Twice a day before the flood, tides in the subterranean water compressed and stretched these pillars. Compressive heating occurred and cracks developed. Just as water circulates through a submerged sponge that is squeezed and stretched, mineral-laden water circulated through cracks in pillars for years before they broke up. Pillar fragments, launched into space by the fountains of the great deep, became meteoroids. In summary, water did it. Tomeoka’s (and Most Evolutionists’) Interpretation. Impacts on an asteroid cracked the rock that was to become this meteorite. Ice was deposited on the asteroid. Impacts melted the ice, allowing liquid water to circulate through the cracks and deposit hundreds of layers of magnesium, calcium, and carbonate bearing minerals. A final impact blasted rocks from this asteroid into space. In summary, impacts did it. Figure 160: Shatter Cone. When a large, crater-forming meteorite strikes the Earth, a shock wave radiates outward from the impact point. The passing shock wave breaks the rock surrounding the crater into meteorite-size fragments having distinctive patterns called shatter cones. (Until shatter cones were associated with impact craters by Robert S. Dietz in 1969, impact craters were often difficult to identify.) If large impacts on asteroids launched asteroid fragments toward Earth as meteorites, a few meteorites should have shatter cone patterns. None have ever been reported. Therefore, meteorites are probably not derived from asteroids. Likewise, impacts have not launched meteorites from Mars. Question 7: Does other evidence support this hypothesis that asteroids and meteoroids came from Earth?Yes. Here are seventeen additional observations that either support the proposed explanation or are inconsistent with other current theories on the origin of asteroids and meteoroids: 1. The materials in meteorites and meteoroids are remarkably similar to those in the Earth’s crust.32 Some meteorites contain very dense elements, such as nickel and iron. Those heavy elements seem compatible only with the denser rocky planets: Mercury, Venus, and Earth—Earth being the densest. A few asteroid densities have been calculated. They are generally low, ranging from 1.2 to 3.3 gm/cm3. The higher densities match those of the Earth’s crust. The lower densities imply the presence of empty space between loosely held rocks or something light such as water ice.33 PREDICTION 35:Rocks in asteroids are typical of the Earth’s crust. Expensive efforts to mine asteroids34 to recover strategic or precious metals will be a waste of money. 2. Meteorites contain different varieties (isotopes) of the chemical element molybdenum, each isotope having a slightly different atomic weight. If, as evolutionists teach, a swirling cloud of gas and dust mixed for millions of years and produced the Sun, its planets, and meteorites, then each meteorite should have about the same combination of these molybdenum isotopes. Because this is not the case,35 meteorites did not come from a swirling dust cloud or any source that mixed for millions of years. 3. Most meteorites36 and some asteroids37 contain metamorphosed minerals, showing that those bodies reached extremely high temperatures, despite a lifetime in the “deep freeze” of outer space. Radioactive decay within such relatively small bodies could not have produced the necessary heating, because too much heat would have escaped from their surfaces. Stranger still, liquid water altered some meteorites38 while they and their parent bodies were heated—sometimes heated multiple times.39 Impacts in space are often proposed to explain this mysterious heating throughout an asteroid or meteroite. However, an impact would raise the temperature only near the point of impact. Before gravel-size fragments from an impact could become uniformly hot, they would radiate their heat into outer space.40 For centuries before the flood, heat was steadily generated within pillars in the subterranean water chamber. As the flood began, the powerful jetting water launched rock fragments into space—fragments of hot, crushed pillars and fragments from the crumbling walls of the ruptured crust. Those rocks became meteoroids and asteroids. 4. Because asteroids came from Earth, they typically spin in the same direction as Earth (counterclockwise, as seen from the North). However, collisions have undoubtedly randomized the spins of many smaller asteroids in the last few thousand years.41 5. Some asteroids have captured one or more moons. [See Figure 156 at top of this page.] Sometimes the “moon” and asteroid are similar in size. Impacts would not create equal-size fragments that could capture each other.42 The only conceivable way for this to happen is if a potential moon enters an asteroid’s expanding sphere of influence while traveling about the same speed and direction as the asteroid. If even a thin gas surrounds the asteroid, the moon will be drawn closer to the asteroid, preventing the moon from being stripped away later. An “exploded planet” would disperse relatively little gas. The “failed planet explanation” meets none of the requirements. The hydroplate theory satisfies all the requirements. Figure 161: Chondrules. The central chondrule above is 2.2 millimeters in diameter, the size of this circle: o. This picture was taken in reflected light. However, meteorites containing chondrules can be thinly sliced and polished, allowing light from below to pass through the thin slice and into the microscope. Such light becomes polarized as it passes through the minerals. The resulting colors identify minerals in and around the chondrules. [Meteorite from Hammada al Hamra Plateau, Libya.] Chondrules (CON-drools) are strange, spherical, BB-size objects found in 86% of all meteorites. To understand the origin of meteorites we must also understand how chondrules formed. Their spherical shape and texture show they were once molten, but to melt chondrules requires temperatures exceeding 3,000°F. How could chondrules get that hot without melting the surrounding rock, which usually has a lower melting temperature? Because chondrules contain volatile substances that would have bubbled out of melted rock, chondrules must have melted and cooled quite rapidly.23 By one estimate, melting occurred in about one-hundredth of a second.24 The standard explanation for chondrules is that small pieces of rock, moving in outer space billions of years ago, before the Sun and Earth formed, suddenly and mysteriously melted. These liquid droplets quickly cooled, solidified, and then were encased inside the rock that now surrounds them. Such vague conditions, hidden behind a veil of space and time, make it nearly impossible to test this explanation in a laboratory. Scientists recognize that this standard story does not explain the rapid melting and cooling of chondrules or how they were encased uniformly in rocks which are radiometrically older than the chondrules.25 As one scientist wrote, “The heat source of chondrule melting remains uncertain. We know from the petrological data that we are looking for a very rapid heating source, but what?”26 Frequently, minerals grade (gradually change) across the boundaries between chondrules and surrounding material.27 This suggests that chondrules melted while encased in rock. If so, the heating sources must have acted briefly and been localized near the center of what are now chondrules. But how could this have happened? The most common mineral in chondrules is olivine.28 Deep rocks contain many BB-size pockets of olivine. Pillars within the subterranean water probably had similar pockets. As the subterranean water escaped from under the crust, pillars had to carry more of the crust’s weight. When olivine reaches a certain level of compression, it suddenly changes into another mineral, called spinel (spin-EL), and shrinks in volume by about 10%.29 (Material surrounding each pocket would not shrink.) Tiny, collapsing pockets of olivine transforming into spinel would generate great heat, for two reasons. First, the transformation is exothermic; that is, it releases heat chemically. Second, it releases heat mechanically, by friction. Here’s why. At the atomic level, each pocket would collapse in many stages—much like falling dominos or the section-by-section crushing of a giant scaffolding holding up an overloaded roof. Within each pocket, as each microscopic crystal slid over adjacent crystals at these extreme pressures, melting would occur along sliding surfaces. The remaining solid structures in the olivine pocket would then carry the entire compressive load—quickly collapsing and melting other parts of the “scaffolding.” The fountains of the great deep expelled pieces of crushed pillars into outer space where they rapidly cooled. Their tumbling action, especially in the weightlessness of space, would have prevented volatiles from bubbling out of the encased liquid pockets within each rock. In summary, chondrules are a by product of the mechanism that produced meteorites—a rapid process that started under the Earth’s crust as the flood began. Also, tidal effects, as described on pages 425–428, limit the lifetime of the moons of asteroids to about 100,000 years.43 This fact and the problems in capturing a moon caused evolutionist astronomers to scoff at early reports that some asteroids have moons. Figure 162: Peanut Asteroids. The fountains of the great deep expelled dirt, rocks, and considerable water from Earth. About half of that water quickly evaporated into the vacuum of space; the remainder froze. Each evaporated gas molecule became an orbiting body in the solar system. Asteroids then formed as explained on pages 298–302. Many are shaped like peanuts. Gas molecules captured by asteroids or released by icy asteroids became their atmospheres. Asteroids with thick atmospheres sometimes captured smaller asteroids as moons. If an atmosphere remained long enough, the moon would lose altitude and gently merge with the low-gravity asteroid, forming a peanut-shaped asteroid. (We see merging when a satellite or spacecraft reenters Earth’s atmosphere, slowly loses altitude, and eventually falls to Earth.) Without an atmosphere, merging becomes almost impossible. Japan’s Hayabusa spacecraft orbited asteroid Itokawa (shown above) for two months in 2005. Scientists studying Itokawa concluded that it consists of two smaller asteroids that merged. Donald Yeomans, a mission scientist and member of NASA’s Jet Propulsion Laboratory, admitted, “It’s a major mystery how two objects each the size of skyscrapers could collide without blowing each other to smithereens. This is especially puzzling in a region of the solar system where gravitational forces would normally involve collision speeds of 2 km/sec.”45 The mystery is easily solved when one understands the role that water played in the origin of comets and asteroids. Notice, a myriad of rounded boulders, some 150 feet in diameter, litter Itokawa’s surface. High velocity water produces rounded boulders; an exploded planet or impacts on asteroids would produce angular rocks. 6. The smaller moons of the giant planets (Jupiter, Saturn, Uranus, and Neptune) are captured asteroids. Most astronomers probably accept this conclusion, but have no idea how these captures could occur.44 As explained earlier in this chapter, for decades to centuries after the flood the radiometer effect, powered by the Sun’s energy, spiraled asteroids outward from Earth’s orbit. Water vapor, around asteroids and in interplanetary space, temporarily thickened asteroid and planet atmospheres. This facilitated aerobraking which allowed massive planets to capture asteroids. Recent discoveries indicate that Saturn’s 313-mile-wide moon, Enceladus (en-SELL-uh-duhs), is a captured asteroid. Geysers at Enceladus’ south pole are expelling water vapor and ice crystals which escape Enceladus and supply Saturn’s E ring.46 That water contains salts resembling Earth’s ocean waters.47 Because asteroids are icy and weak, they would experience strong tides if captured by a giant planet. Strong tides would have recently48 generated considerable internal heat, slowed the moon’s spin, melted ice, and boiled deep reservoirs of water. Enceladus’ spin has almost stopped, its internal water is being launched (some so hot that it becomes a plasma),49 and its surface near the geysers has buckled, probably due to the loss of internal water. Because the material for asteroids and their organic matter came recently from Earth, water is still jetting from cold Enceladus’ surprisingly warm south pole, and “dark green organic material”50 is on its surface. 7. A few asteroids suddenly develop comet tails, so they are considered both asteroid and comet. The hydroplate theory says that asteroids are weakly joined piles of rocks and ice. If such a pile cracked slightly, perhaps due to an impact by space debris, then internal ice, suddenly exposed to the vacuum of space, would violently vent water vapor and produce a comet tail. The hydroplate theory explains why comets are so similar to asteroids. 8. A few comets have nearly circular orbits within the asteroid belt. Their tails lengthen as they approach perihelion and recede as they approach aphelion. If comets formed beyond the planet Neptune, it is highly improbable that they could end up in nearly circular orbits in the asteroid belt.51 So, these comets almost certainly did not form in the outer solar system. Also, comet ice that near the Sun would evaporate relatively quickly. Only the hydroplate theory explains how comets (icy rock piles) recently entered the asteroid belt. 9. If asteroids passing near Earth came from the asteroid belt, too many of them have diameters less than 50 meters,52 and too many have circular orbits.53 However, we would expect this if the rocks that formed asteroids were launched from Earth. 10. Computer simulations, both forward and backward in time, show that asteroids traveling near Earth have a maximum expected lifetime of only about a million years. They “quickly” collide with the Sun.54 This raises doubts that all asteroids began 4,600,000,000 years ago as evolutionists claim—living 4,600 times longer than the expected lifetime of near-Earth asteroids. 11. Earth has one big moon and several small moons—up to 650 feet in diameter.55 The easiest explanation for the small moons is that they were launched from Earth with barely enough velocity to escape Earth’s gravity. (To understand why the largest of these small moons is about 650 feet in diameter, see Endnote 8.) 12. Asteroids 3753 Cruithne and 2000 AA29 are traveling companions of Earth.56 They delicately oscillate, in a horseshoe pattern, around two points that lie 60° (as viewed from the Sun) forward and 60° behind the Earth but on Earth’s nearly circular orbit. These points, predicted by Lagrange in 1764 and called Lagrange points, are stable places where an object would not move relative to the Earth and Sun if it could once occupy either point going at zero velocity relative to the Earth and Sun. But how could a slowly moving object ever reach, or get near, either point? Most likely, it barely escaped from Earth. Also, Asteroid 3753 could not have been in its present orbit for long, because it is so easy for a passing gravitational body to perturb it out of its stable niche. Time permitting, Venus will pass near this asteroid 8,000 years from now and may dislodge it.57 13. Furthermore, Jupiter has two Lagrange points on its nearly circular orbit. The first, called L4, lies 60° (as seen from the Sun) in the direction of Jupiter’s motion. The second, called L5, lies 60° behind Jupiter. Visualize planets and asteroids as large and small marbles rolling in orbitlike paths around the Sun on a large frictionless table. At each Lagrange point is a bowl-shaped depression that moves along with each planet. Because there is no friction, small marbles (asteroids) that roll down into a bowl normally pick up enough speed to roll back out. However, if a chance gravitational encounter slowed one marble right after it entered a bowl, it might not exit the bowl. Marbles trapped in a bowl would normally stay 60° ahead of or behind their planet, gently rolling around near the bottom of their moving bowl. One might think an asteroid is just as likely to get trapped in Jupiter’s leading bowl as its trailing bowl—a 50–50 chance, as with the flip of a coin. Surprisingly, 1068 asteroids are in Jupiter’s leading (L4) bowl, but only 681 are in the trailing bowl.69 This shouldn’t happen in a trillion trials if an asteroid is just as likely to get trapped at L4 as L5. What concentrated so many asteroids near the L4 Lagrange point? According to the hydroplate theory, asteroids formed near Earth’s orbit. Then, the radiometer effect spiraled them outward, toward the orbits of Mars and Jupiter. Some spiraled through Jupiter’s circular orbit and passed near both L4 and L5. Jupiter’s huge gravity would have slowed those asteroids that were moving away from Jupiter but toward L4. That braking action would have helped some asteroids settle into the L4 bowl. Conversely, asteroids that entered L5 were accelerated toward Jupiter, so they would quickly be pulled out of L5 by Jupiter’s gravity. The surprising excess of asteroids near Jupiter’s L4 is what we would expect based on the hydroplate theory. Figure 163: Asteroid Belt and Jupiter’s L4 and L5. The size of the Sun, planets, and especially asteroids are magnified, but their relative positions are accurate. About 90% of the 30,000 precisely known asteroids lie between the orbits of Mars and Jupiter, a doughnut-shaped region called the asteroid belt. A few small asteroids cross Earth’s orbit. Jupiter’s Lagrange points, L4 and L5, lie 60° ahead and 60° behind Jupiter, respectively. They move about the Sun at the same velocity as Jupiter, as if they were fixed at the corners of the two equilateral triangles shown. Items 12 and 13 explain why so many asteroids have settled near L4 and L5, and why significantly more oscillate around L4 than L5. 14. Without the hydroplate theory, one has difficulty imagining situations in which an asteroid would (a) settle into one of Jupiter’s Lagrange points, (b) capture a moon, especially a moon with about the same mass as the asteroid, or (c) have a circular orbit, along with its moon, about their common center of mass. If all three happened to an asteroid, astronomers would be shocked; no astronomer would have predicted that it could happen to a comet. Nevertheless, an “asteroid” discovered earlier, named 617 Patroclus, satisfies (a)–(c). Patroclus and its moon, Menoetius, have such low densities that they would float in water; therefore, both are probably comets70—dirty, fluffy snowballs. Paragraphs 5, 7, 8, and 13 (above) explain why these observations make perfect sense with the hydroplate theory. 15. As explained in “Shallow Meteorites,” meteorites are almost always found surprisingly near Earth’s surface. The one known exception is in southern Sweden, where 40 meteorites and thousands of grain-size fragments of one particular type of meteorite have been found at different depths in a few limestone quarries. The standard explanation is that all these meteorites somehow struck this same small area over a 1–2-million-year period about 480 million years ago.71 A more likely explanation is that some meteorites, not launched with enough velocity to escape Earth during the flood, fell back to Earth. One or more meteorites fragmented on reentering Earth’s atmosphere. The pieces landed in mushy, recently-deposited limestone layers in southern Sweden. 16. Light spectra (detailed color patterns, much like a long bar code) from certain asteroids in the outer asteroid belt imply the presence of organic compounds, especially kerogen, a coal-tar residue.72 No doubt the kerogen came from plant life. Life as we know it could not survive in such a cold region of space, but common organic matter launched from Earth could have been preserved. 17. Many asteroids are reddish and have light characteristics showing the presence of iron.73 On Earth, reddish rocks almost always imply iron oxidized (rusted) by oxygen gas. Today, oxygen is rare in outer space. If iron on asteroids is oxidized, what was the source of the oxygen? Answer: Water molecules, surrounding and impacting asteroids, dissociated (broke apart), releasing oxygen. That oxygen then combined chemically with iron on the asteroid’s surface, giving the reddish color. Mars, often called the red planet, derives its red color from oxidized iron. Again, oxygen contained in water vapor launched from Earth during the flood, probably accounts for Mars’ red color. Mars’ topsoil is richer in iron and magnesium than Martian rocks beneath the surface. The dusty surface of Mars also contains carbonates, such as limestone.74 Because meteorites and Earth’s subterranean water contained considerable iron, magnesium, and carbonates, it appears that Mars was heavily bombarded by meteorites and water launched from Earth’s subterranean chamber. [See “The Origin of Limestone” on pages 224–229.] Those who believe that meteorites came from asteroids have wondered why meteorites do not have the red color of most asteroids.75 The answer is twofold: (a) as explained on page 301, meteorites did not come from asteroids but both came from Earth, and (b) asteroids contain oxidized iron, as explained above, but meteorites are too small to attract an atmosphere gravitationally. Figure 164: Salt of the Earth. On 22 March 1998, this 2 3/4 pound meteorite landed 40 feet from boys playing basketball in Monahans, Texas. While the rock was still warm, police were called. Hours later, NASA scientists cracked the meteorite open in a clean-room laboratory, eliminating any possibility of contamination. Inside were salt (NaCl) crystals 0.1 inch (3 mm) in diameter and liquid water!58 Some of these salt crystals are shown in the blue circle, highly magnified and in true color. Bubble (B) is inside a liquid, which itself is inside a salt crystal. Eleven quivering bubbles were found in about 40 fluid pockets. Shown in the green circle is another bubble (V) inside a liquid (L). The length of the horizontal black bar represents 0.005 mm, about 1/25 the diameter of a human hair. NASA scientists who investigated this meteorite believe that it came from an asteroid, but that is highly unlikely. Asteroids, having little gravity and being in the vacuum of space, cannot sustain liquid water, which is required to form salt crystals. (Earth is the only planet, indeed the only body in the solar system, that can sustain liquid water on its surface.) Nor could surface water (gas, liquid, or solid) on asteroids withstand high-velocity impacts. Even more perplexing for the evolutionist: What is the salt’s origin? Also, what accounts for the meteorite’s other contents: potassium, magnesium, iron, and calcium—elements abundant on Earth, but as far as we know, not beyond Earth?59 Dust-sized meteoroids often come from comets. Most larger meteoroids are rock fragments that never merged into a comet or asteroid. Much evidence supports Earth as the origin of meteorites. - Minerals and isotopes in meteorites are remarkably similar to those on Earth.32 - Some meteorites contain sugars,60 salt crystals containing liquid water,61 and possible cellulose.62 - Other meteorites contain limestone,63 which, on Earth, forms only in liquid water. - Three meteorites contain excess amounts of left-handed amino acids64—a sign of once-living matter. - A few meteorites show that “salt-rich fluids analogous to terrestrial brines” flowed through their veins.65 - Some meteorites have about twice the heavy hydrogen concentration as Earth’s water today.66 As explained in the preceding chapter and in “Energy in the Subterranean Water” here, this heavy hydrogen came from the subterranean chambers. - About 86% of all meteorites contain chondrules, which are best explained by the hydroplate theory. - Seventy-eight types of living bacteria have been found in two meteorites after extreme precautions were taken to avoid contamination.67 Bacteria need liquid water to live, grow, and reproduce. Obviously, liquid water does not exist inside meteoroids whose temperatures in outer space are near absolute zero (-460°F). Therefore, the bacteria must have been living in the presence of liquid water before being launched into space. Once in space, they quickly froze and became dormant. Had bacteria originated in outer space, what would they have eaten? Water on MarsWater recently and briefly flowed at various locations on Mars.76 Photographic comparisons show that some water flowed within the last 2–5 years!77 Water is now stored as ice at Mars’ poles78 and in surface soil. Mars’ stream beds usually originate on crater walls rather than in ever smaller tributaries as on Earth.79 Rain formed other channels.80 Martian drainage channels and layered strata are found at almost isolated 200 locations.81 Most gullies are on crater slopes at high latitudes82—extremely cold slopes that receive little sunlight. One set of erosion gullies is on the central peak of an impact crater!83 Figure 165: Erosion Channels on Mars. These channels frequently originate in scooped-out regions, called amphitheaters, high on a crater wall. On Earth, where water falls as rain, erosion channels begin with narrow tributaries that merge with larger tributaries and finally, rivers. Could impacts of comets or icy asteroids have formed these craters, gouged out amphitheaters, and melted the ice—each within seconds? Mars, which is much colder than Antarctica in the winter, would need a heating source, such as impacts, to produce liquid water. Today, Mars is cold, averaging -80°F (112 Fahrenheit degrees below freezing). Water on Mars should be ice, not liquid water. Mars’ low atmospheric pressures would hasten freezing even more.84 Water probably came from above. Soon after Earth’s global flood, the radiometer effect caused asteroids to spiral out to the asteroid belt, just beyond Mars. This gave asteroids frequent opportunities to collide with Mars. When crater-forming impacts occurred, large amounts of debris were thrown into Mars’ atmosphere. Mars’ thin atmosphere and low gravity allowed the debris to settle back to the surface in vast layers of thin sheets—strata. PREDICTION 36Most sediments taken from layered strata on Mars and returned to Earth will show that they were deposited through Mars’ atmosphere, not through water. (Under a microscope, water deposited grains have nicks and gouges, showing that they received many blows as they tumbled along stream bottoms. Sediments deposited through an atmosphere receive few nicks.) Impact energy (and heat) from icy asteroids and comets bombarding Mars released liquid water, which often pooled inside craters or flowed downhill and eroded the planet’s surface.87 (Most liquid water soaked into the soil and froze.) Each impact was like the bursting of a large dam here on Earth. Brief periods of intense, hot rain and localized flash floods followed.88 These Martian hydrodynamic cycles quickly “ran out of steam,” because Mars receives relatively little heat from the Sun. While the consequences were large for Mars, the total water was small by Earth’s standards—about twice the water in Lake Michigan. Today, when meteorites strike icy soil on Mars, some of that ice melts. When this happens on a crater wall, liquid water flows down the crater wall, leaving the telltale gullies that have shocked the scientific community.77 PREDICTION 37As has been discovered on the Moon and apparently on Mercury, frost will be found within asteroids and in permanently shadowed craters on Mars. This frost will be rich in heavy hydrogen. Are Some Meteorites from Mars?Widely publicized claims have been made that at least 30 meteorites from Mars have been found. With international media coverage in 1996, a few scientists also proposed that one of these meteorites, named ALH84001, contained fossils of primitive life. Later study rejected that claim. The wormy-looking shapes discovered in a meteorite [supposedly] from Mars turned out to be purely mineralogical and never were alive.89 The 30 meteorites are presumed to have come from the same place, because they contain similar ratios of three types of oxygen: oxygen weighing 16, 17, and 18 atomic mass units. (That presumption is not necessarily true, is it?) A chemical argument then indirectly links one of those meteorites to Mars, but the link is more tenuous than most realize.90 That single meteorite had tiny glass nodules containing dissolved gases. A few of these gases (basically the noble gases: argon, krypton, neon, and xenon) had the same relative abundances as those found in Mars’ atmosphere in 1976. (Actually, a later discovery shows that the mineralogy of these meteorites differs from that of almost all Martian rock.91) Besides, if two things are similar, it does not mean that one came from the other. Similarity in the relative abundances of the noble gases in Mars’ atmosphere and in one meteorite may be because those gases originated in Earth’s preflood subterranean chamber. Rocks and water from the subterranean chamber may have transported those gases to Mars. Could those 30 meteorites have come from Mars? To escape the gravity of Mars requires a launch velocity of 3 miles per second. Additional velocity is then needed to transfer to an orbit intersecting Earth, 34–236 million miles away. Supposedly, one or more asteroids slammed into Mars and blasted off millions of meteoroids. Millions are needed, because less than one in a million92 would ever hit Earth, be large enough to survive reentry, be found, be turned over to scientists, and be analyzed in detail. Besides, if meteorites can come to Earth from Mars, many more should have come from the Moon—but haven’t.93 For an impact suddenly to accelerate, in a fraction of a second, any solid from rest to a velocity of 3 miles per second requires such extreme shock pressures that much of the material would melt, if not vaporize.94 All 30 meteorites should at least show shock effects. Some do not. Also, Mars should have at least six giant craters if such powerful blasts occurred, because six different launch dates are needed to explain the six age groupings the meteorites fall into (based on evolutionary dating methods). Such craters are hard to find, and large, recent impacts on Mars should have been rare. Then there are energy questions. Almost all impact energy is lost as shock waves and ultimately as heat. Little energy remains to lift rocks off Mars. Even with enough energy, the fragments must be large enough to pass through Mars’ atmosphere. To see the difficulty, imagine throwing a ball high into the air. Then visualize how hard it would be to throw a handful of dust that high. Atmospheric drag, even in Mars’ thin atmosphere, absorbs too much of the smaller particles’ kinetic energy. Finally, for large particles to escape Mars, the expelling forces must be focused, as occurs in a gun barrel or rocket nozzle. For best results, this should be aimed straight up, to minimize the path length through the atmosphere. A desire to believe in life on Mars produced a type of “Martian mythology” that continues today. In 1877, Italian astronomer Giovanni Schiaparelli reported seeing grooves on Mars. The Italian word for groove is “canali”; therefore, many of us grew up hearing about “canals” on Mars—a mistranslation. Because canals are man-made structures, people started thinking about “little green men” on Mars. In 1894, Percival Lowell, a wealthy, amateur astronomer with a vivid imagination, built Lowell Observatory primarily to study Mars. Lowell published a map showing and naming Martian canals, and wrote several books: Mars (1895), Mars and Its Canals (1906), and Mars As the Abode of Life (1908). Even into the 1960s, textbooks displayed his map, described vegetative cycles on Mars, and explained how Martians may use canals to convey water from the polar ice caps to their parched cities. Few scientists publicly disagreed with the myth, even after 1949 when excellent pictures from the 200-inch telescope on Mount Palomar were available. Those of us in school before 1960 were directly influenced by such myths; almost everyone has been indirectly influenced. Artists, science fiction writers, and Hollywood helped fuel this “Martian mania.” In 1898, H. G. Wells wrote The War of the Worlds telling of strange-looking Martians invading Earth. In 1938, Orson Welles, in a famous radio broadcast, panicked many Americans into thinking New Jersey was being invaded by Martians. In 1975, two Viking spacecraft were sent to Mars to look for life. Carl Sagan announced, shortly before the tests were completed, that he was certain life would be discovered—a reasonable conclusion, if life evolved. The prediction failed. In 1996, United States President Clinton read to a global television audience, “More than 4 billion years ago this piece of rock [ALH84001] was formed as a part of the original crust of Mars. After billions of years, it broke from the surface and began a 16-million-year journey through space that would end here on Earth.” “... broke from the surface ...”? The myth is still alive. Final ThoughtsAs with the 24 other major features listed on page 106 [of the book, In the Beginning], we have examined the origin of asteroids and meteoroids from two directions: “cause-to-effect” and “effect-to-cause.” Cause-to-Effect. We saw that given the assumption listed on page 115 [of the book, In the Beginning], consequences naturally followed: subterranean water became supercritical, the fountains of the great deep erupted; large rocks, muddy water, and water vapor were launched into space; gas and gravity assembled asteroids; and gas pressure powered by the Sun’s energy (the radiometer effect) herded asteroids into the asteroid belt. Isolated rocks still moving in the solar system are meteoroids. Effect-to-Cause. We considered seventeen effects (pages 302–306)[of the book, In the Beginning], each incompatible with present theories on the origin of asteroids and meteoroids. Each effect was evidence that many rocks and large volumes of water vapor were launched from Earth. Portions of Part III will examine this global flood from a third direction: historical records from claimed eyewitnesses. All three perspectives reinforce each other, illuminating in different ways this catastrophic event. To access the footnotes for this article, click here.

http://4thdayalliance.com/articles/solar-system/origin-of-asteroids/

Applying Doppler Effect to Moving Galaxies Overview: Students make the observation that farther galaxies move away faster, and check that a model of an expanding universe makes predictions that match with those observations. Physical resources: Expanding universe model Electronic resources: Virtual spectroscopy Observations of moving galaxies: - Motivating question: How can we use the idea of redshift to figure out the velocity of objects? Students brainstorm ideas with group. - How do we know what was emitted? Introduce spectral lines as the photon we know must have been emitted with a certain energy, in our case, we'll look at line emission from hydrogen atoms. - Virtual spectroscope activity: (MiniSpectroscopy) - Examine Hydrogen spectrum at rest, predict how "example galaxy" is moving, relative to Earth. (Peak is at a longer wavelength, so it is moving away from us.) - Give students only the spectra of galaxies A through D - What direction are they moving? (away from Earth, because peak of emission is at a longer wavelength than it is when hydrogen is at rest) - Challenge: put them in order by the speed (slowest to fastest) they are moving away from Earth. - Now, give students the images of galaxies A through D - What's different about these galaxies? (angular diameter) - Given that most galaxies are about the same linear diameter, put them in order by their distance from Earth, closest to furthest. - Students should describe the pattern in these observations, and put their description on the whiteboard. (The order is the same, further galaxies move away faster.) - Instructor introduces Hubble's law as a restatement of this observation: Galaxies that are farther away move away from us faster. Model of expanding universe, to explain Hubble's law observations above: - Introduce two-dimensional "expanding universe" model which we've taken an image of at two different times - Label "smaller" universe as time t = 0, and "larger" universe as time t = 10 seconds. - Have groups of students "live" in galaxy A, B or C and have them make predictions of the following for each of the two other labeled galaxies, as well as another galaxy of their choice: - Distance from your galaxy to other galaxy at time t = 0 seconds (cm) - Distance from your galaxy to other galaxy at time t = 10 seconds (cm) - Change in distance (cm) - Change in time (sec, all should be 10 seconds) - Speed = change in distance / change in time (cm / sec) - Direction of motion (description, or arrow) - Have students populate classroom prediction table - Summarize important patterns seen in predictions: Galaxies at a greater distance move faster. - Have students line up their "home galaxy" while holding both "universes" up to the light, and describe what has happened to all the other galaxies (they have moved away from the home galaxy, on a line connecting the home galaxy to the other galaxy.) Then have them switch their "Home galaxy" to the other two labeled galaxies, in turn. (All galaxies will see this pattern of all others moving away). - Refined prediction: Galaxies at a greater distance move faster, and move away from each other along a line connecting the two. From any galaxy, all others look like they are moving away. - These predictions match up with the observations we've made about actual galaxies in our universe, so we can't rule out the "expanding universe" model. - Some students have difficulty identifying what information they should extract from the spectra when comparing the sample of galaxies. They may think the intensity of the peak is what they should order by, instead of the location of the peak on the wavelength (energy) scale. - Many students have difficulty separating the observations from the models in this activity. If so, clarify with the assessment question below. - Can we determine redshifts for galaxies that do not have emission lines? (no, we must know the energy at which the photons were originally emitted). - What if we observed every galaxy moving toward us, with further galaxies moving toward us faster? How would that change our model to explain the observations? (contracting universe). - Which is a statement of the Doppler effect, and which is a statement of Hubble's Law? - When we observe galaxies moving away from us, we receive lower energy photons compared to what that galaxy actually emits. (Doppler effect) - Galaxies moving away from us faster are also further away (Hubble's Law). - Galaxies moving towards us give us photons that are higher energy than when they were emitted (Doppler effect). - Closer galaxies are moving away from us slower (Hubble's Law). - Discuss what each deals with: Hubble's law relates speed to distance, and Doppler effect relates change in energy of photons to speed of motion. - Image of review page of notes: (Hubble's law 2) < return to Investigation 6

http://ocw.mit.edu/high-school/courses/chandra-astrophysics-institute/investigations/investigation-6/activity-3/

Get answers to your child's growth, nutrition, and feeding behavior questions. School-age child development describes the expected physical, emotional, and mental abilities of children ages 6 - 12. School-age children usually have smooth and strong motor skills. However, their coordination (especially eye-hand), endurance, balance, and physical abilities vary. Fine motor skills may also vary widely. These skills can affect a child's ability to write neatly, dress appropriately, and perform certain chores, such as making beds or doing dishes. There will be big differences in height, weight, and build among children of this age range. It is important to remember that genetic background, as well as nutrition and exercise, may affect a child's growth. There can also be a big difference in the age at which children begin to develop secondary sexual characteristics. For girls, secondary sex characteristics include: For boys, they include: Early school-age children should be able to use simple, but complete sentences that average five to seven words. As the child goes through the elementary school years, grammar and pronunciation become normal. Children use more complex sentences as they grow. Language delays may be due to hearing or intelligence problems. In addition, children who are unable to express themselves well may be more likely to have aggressive behavior or temper tantrums. A 6-year-old child normally can follow a series of three commands in a row. By age 10, most children can follow five commands in a row. Children who have a problem in this area may try to cover it up with backtalk or clowning around. They will rarely ask for help because they are afraid of being teased. Frequent physical complaints (such as sore throats, tummy aches, arm or leg pain) may simply be due to a child's increased body awareness. Although there is often no physical evidence for such complaints, the complaints should be investigated to rule out possible health conditions, and to assure the child that the parent is concerned about his or her well-being. Peer acceptance becomes more important during the school-age years. Children may take part in certain behaviors to be part of "the group." Talking about these behaviors with your child will allow the child to feel accepted in the group, without crossing the boundaries of the family's behavior standards. Friendships at this age tend to be mainly with members of the same sex. In fact, younger school-age children often talk about members of the opposite sex as being "strange" or "awful." Children become less negative about the opposite sex as they get closer to adolescence. Lying, cheating, and stealing are all examples of behaviors that school-age children may "try on" as they learn how to negotiate the expectations and rules placed on them by family, friends, school, and society. Parents should deal with these behaviors privately (so that the child's friends don't tease them). Parents should show forgiveness, and punish in a way that is related to the behavior. An ability to pay attention is important for success both at school and at home. A 6-year-old should be able to focus on a task for at least 15 minutes. By age 9, a child should be able to focus attention for about an hour. It is important for the child to learn how to deal with failure or frustration without losing self-esteem. Safety is important for school-age children. Feigelman S. Middle childhood. In: Kliegman RM, Behrman RE, Jenson HB, Stanton BF, eds. Nelson Textbook of Pediatrics. 18th ed. Philadelphia, Pa: Saunders Elsevier; 2007:chap 11. © 2011 University of Maryland Medical Center (UMMC). All rights reserved. UMMC is a member of the University of Maryland Medical System, 22 S. Greene Street, Baltimore, MD 21201. TDD: 1-800-735-2258 or 1.866.408.6885

http://www.umm.edu/ency/article/002017.htm

GMAT Coordinate Geometry August 1, 2012 The key to many GMAT coordinate geometry questions is to remember that coordinate geometry is just another way of expressing the possible solutions to a two variable equation. Each point on the line in a coordinate plane corresponds to a solution for the equation of that line. The base equation for a line is y = mx + b, where b is the y intercept, or the point at which the line crosses the y-axis, and m is the slope, or the steepness of the line. More specifically, the slope of a line is the change in the y coordinates divided by the change in the x coordinates between any two points on the line. While understanding the basic format for an equation of a line can be very useful on the GMAT quantitative section, you will encounter GMAT problems in which it is faster and easier to think of the problem in algebraic terms. In such cases you should think of the equation as an algorithm that will produce the y value given any x value. This is the reason that the x values are sometimes referred to as inputs and the y values as outputs. For example, if your answer choices are solution sets and you are asked to determine which option is on the line given in the y = mx + b form, rather than graphing the line and trying to determine which point falls on it, which is especially difficult as you will not have graph paper, you can plug each x value into the equation and determine which one produces the appropriate y value. On test day, the key is to remember that coordinate geometry is just a way of expressing algebraic concepts visually. Thus, we can often treat these problems as algebra rather than as geometry. To see this in action, try the problem below. In the xy-coordinate system, if (m, n) and (m 1 2, n 1 k) are two points on the line with the equation x 5 2y 1 5, then k 5 Step 1: Analyze the Question For any question involving the equation of a line, a good place to start is the slope-intercept form of the line, y = mx 1 b. Remember that if you have two points on a line, you can derive the entire equation, and if you have an equation of the line, you can calculate any points on that Step 2: State the Task We are solving for k, which is the amount by which the y-coordinate increases when the x-coordinate increases Step 3: Approach Strategically The slope of a line is the ratio between the change in y and the change in x. In other words, every time the x-coordinate increases by 1, the y-coordinate increases by the amount of the slope. The equation of the line in the question stem is defined as x = 2y + 5. We must isolate y to have slope-intercept form: So the slope of this line is 1/2 . This means that for every change of +1 in the x direction, there is a change of + 1/2 in the y direction. Then we know that, because there is an increase in 2 units in the x direction when moving from m to m + 2, there must be a change of 1 unit in the y direction when moving from n to n + k. So k = 1. Since there are variables that eventually cancel (m and n are not part of the answers), we can Pick Numbers. Let’s say that you choose the y-coordinate of the point (m, n) to be 0 to allow for easier calculations. Using the equation we’re given to relate x- and y-coordinates, we can calculate So (m, n) is the point (5, 0). Now we’ll plug our values of m and n into the next point: (m + 2, n + k). That yields (7, k). All we have to do is plug an x-coordinate of 7 into the equation to solve for k, the

http://blog.kaplangmat.com/2012/08/01/gmat-coordinate-geometry/

Early federal bankruptcy laws were typically temporary responses to bad economic conditions, created to punish debtors. Modern laws emphasize rehabilitating distressed debtors. Pre-Revolutionary U.S. bankruptcy laws, by and large, resembled English laws. which were creditor-oriented, punitive laws created in the late 1500s and 1600s. Imprisonment for debt was permitted in every colony and later every state. In fact, “even as late as 1785, the Pennsylvania Bankruptcy Act allowed flogging for ‘convicted bankrupts’ ” (Levin and Rogers, July/August 2007). In the early 19th century, attitudes toward debtors began to slowly transform. The public became more sympathetic toward debtors, possibly realizing that harsh punishment often did not provide material compensation to creditors. The law began protecting debtors to encourage debtor-creditor resolutions, nullifying debts of cooperative debtors trying to reduce their outstanding monetary obligations.

http://www.pluris.com/bankruptcy-more-information

Identifying Similarities and Differences From Web 2.0 That Works: Marzano & Web 2.0 Enhance students' understanding of and ability to use knowledge by engaging them in mental processes that involve identifying ways items are alike and different. Generalizations From Research - Presenting students with explicit guidance in identifying similarities and differences enhances students’ understanding of and ability to use knowledge. - Asking students to independently identify similarities and differences enhances students’ understanding of and ability to use knowledge. - Representing similarities and differences in graphic or symbolic form enhances students’ understanding of and ability to use knowledge. - Identification of similarities and differences can be accomplished in a variety of ways. The identification of similarities and differences is a highly robust activity. - Venn Diagrams - Comparison Matrix - Analogies Organizer Recommendations & Ideas - Use comparing, classifying, metaphors, and analogies when having students compare similarities and differences. - Give students a model of the steps for engaging in the process. - Use a familiar context to teach students these steps. - Use Graphic Organizers as a visual tool to represent similarities and differences. - Guide students as they engage in this process. Gradually give less structure and less guidance. Information presented above in the definition is from McREL, and generalizations from research, and recommendations & ideas is from Classroom Instruction that Works: Research-Based Strategies for Increasing Student Achievement (ASCD)(Robert J. Marzano, Debra J. Pickering, Jane E. Pollock). Web 2.0 Connections "Y" Under each category indicates that this tool can be used with this strategy. "Free +" Indicates that the tool is free at the basic level, but that more advanced versions are available at a cost. SD = Identifying Similarities and Differences CL = Cooperative Learning SNT = Summarizing and Note-Taking ER = Reinforcing Effort and Providing Recognition HP = Homework and Practice NR = Nonlinguistic Representation OF = Setting Objectives and Providing Feedback HYP = Generating and Testing Hypotheses QCO = Questions, Cues, and Advance Organizers |VoiceThread||[]||Online Collaborative Slideshows||Free+||Y||Y||Y||Y||Y||his online tool allows for collaboration on projects from any internet-connected computer. Products can be very rich visually and aurally, and the commenting feature is excellent. Comments can be made through text, recording on the computer, or calling in via phone (fees apply).| |Exploratree||[]||Online Graphic Organizer||Free||Y||Y||Y||Y||Y||Y||Interactive, pre-made graphic organizers that can be edited online| |MediaWiki||[]||Open Source wiki software||Free||Y||Y||Y||Y||Y||Y||Y||Y||Y||Wiki engine that powers [http://www.wikipedia.org Wikipedia] as well as this wiki.| |iGoogle||[]||Online Personal Homepage||Free||Y||Y||Y||Y||Y||Y||Y||Y||Integrates with all Google tools| |Netvibes||[]||Online Personal Homepage||Free||Y||Y||Y||Y||Y||Y||Y||Y| |Backpack||[]||Online Personal Organizer||Free +||Y||Y||Y||Y||Y||Y||Y| |Ning||[]||Community Websites||Free||Y||Y||Y||Y||Y||Y||Y||Y||Y||Can be used for a classroom website| |EditGrid||[]||Online Spreadsheets||Free +||Y||Y||Y||Y||Y||Integrates with Facebook and iPhone| |Diigo||[]||Online Social Annotation||Free||Y||Y||Y||Y||Y||Y| |Writeboard||[]||Online Collaborative Note-Taking||Free||Y||Y||Y||Y||Y||Y||Y||Integrates with BackPack| |stu.dicio.us||[]||Online Note-Taking||Free||Y||Y||Y||Y||Y||Y||Y||Currently under redevelopment| |pbWiki||[]||Wiki Hosting Website||Free||Y||Y||Y||Y||Y||Y||Y||Y||Y| |YackPack||[]||Online voice Messaging||Free||Y||Y||Y||Y||Y||Y||Links directs to YackPack for Educators information| |Ajax13||[]||Online Graphic Editor||Free||Y||Y||Y||Y||Y||Y||Y||Requires Firefox 1.5 (or higher) Browser| |SketchUp||[]||3D Drawing Software||Free||Y||Y||Y||Y||Integrates with Google Earth| |Thinkature||[]||Online Collaboration & Whiteboard||Free||Y||Y||Y||Y||Y||Y||Y||Y| |ScanR||[]||Online Scanner, Copier, Fax from Phone or Digital Camera||Free +||Y||Y||Y||Y||Y||Y| |Flickr||[]||Photo Hosting Website||Free +||Y||Y||Y||Y||Part of Zoho Suite of Online Apps| |Zoho Wiki||[]||Wiki Hosting Website||Free||Y||Y||Y||Y||Y||Y||Y||Y||Y||Part of Zoho Suite of Online Apps| |Zoho Sheet||[]||Online Spreadsheets||Free||Y||Y||Y||Y||Y||Part of Zoho Suite of Online Apps| |Google SpreadSheets||[]||Online Spreadsheets||Free||Y||Y||Y||Y||Y||Also contains Documents & Presentations| |WordPress||[]||Blog Hosting Website||Free||Y||Y||Y||Y||Y||Y| |Technorati||[]||Blog Search Engine||Free||Y||Y||Y||Y||Y||Y| |EduBlogs||[]||Blog Hosting Website||Free||Y||Y||Y||Y||Y||Y| |Blogger||[]||Blog Hosting Website||Free||Y||Y||Y||Y||Y||Y| |Wikispaces||[]||Wiki Hosting Website||Free||Y||Y||Y||Y||Y||Y||Y||Y||Y| Examples from teachers, students, classrooms, schools Please click on the "edit" tab to share your examples here. Contributor: Louise Maine In order to understand the classification of animals in the phyla and some classes of phyla, students were given two organisms each of two particular classes or phyla to determine similarities and differences among them. Information was collected from the class and put together in a large document to show these in relation to the animal kingdom. As a class, the information was discussed to show how the similarities and differences among these animals placed them in a particular class and phylum. Contributor: Tiffany Zarling In order to simplify the voice recording process for book trailers, students use www.voki.com. They create a free account, create their own avatar, and can record their voices in a variety of ways (text to speech, phone call, uploading, or recording directly). It IS a free+ tool. The free accounts record up to one minute of voice. You can email, embed, or house your finished Voki through the website as well. It was a great motivating tool for senior level English students!

http://web2thatworks.com/index.php?title=Identifying_Similarities_and_Differences

Northern Prairie Wildlife Research Center A moth is the sexually mature adult life stage and serves three main functions in the life cycle: mating, dispersal, and oviposition. Many moths feed on nectar or a liquid sugar source, which primarily serves as fuel for flight. Some species of macromoths do not have functional mouthparts and cannot feed, therefore they are relatively short-lived and in turn will exhibit a short flight period. Dispersal and flight activity. Typically moths possess two pairs of wings, a pair of forewings and a pair of hindwings. The forewings are attached to the second thoracic segment, the mesothorax, while the hindwings are attached to the third thoracic segment, the metathorax. Moths capable of flight, which is the primary means for local and long distance movement, may beat their wings up to 60 or 70 times per second. However, not all moths have wings, and not all moths with wings can fly. Individuals that do not have wings, as well as those that have wings but are flightless, do not have flight muscles. Typically the female is subject to the loss of flight muscles, which is accompanied by a higher capacity for egg production. The male of the species has fully developed wings and is flight capable. Examples of species with wingless females are the lymantriids Orgyia antiqua and Orgyia pseudotsugata, and the geometrids Erannis tiliaria, Operopthera bruceata, and Operopthera danbyi. The period for flight may be characteristic for a species and must be assessed in two types of time periods, the daily rhythm and the seasonal pattern. The majority of macromoths are night-flying species while a minority of species fly during the day. Some of the day-flying macromoths exhibit highly contrasting colors on their wings as exemplified by many of the arctiids such as Gnophaela vermiculata, Tyria jacobaeae, Leptarctia californiae, and Platyprepia virginalis; the male of the lymantriid Orgyia antiqua; the geometrids Rheumaptera subhastata and Mesoleuca gratulata; the noctuids Alypia langtoni and Schinia walsinghami; the saturniids Hemileuca eglanterina and Saturnia mendocino; and the sphingids Hemaris diffinis and Proserpinus clarkiae. The behavior of being either a night or day flying moth is characteristic to most species, however, a few of the night flying species, such as Hyles lineata, may be seen on the wing during the day. The time of season and the length of time for the flight period of a species may also exhibit a diagnostic pattern. Most species occur at certain times of the year and may be present for a period of three to six weeks. For instance, the arctiid Lophocampa argentata will be in flight during the last few days of July and the first three weeks of August with a peak in flight around the end of the first week in August. Similarly, the males of the geometrids Operopthera bruceata and Operopthera danbyi are only in flight from the middle of November to the last week of December. On the other hand a few species may have individuals in flight throughout much of the year. For instance, the geometrids Orthonama centrostrigaria and Sabulodes aegrotata fly from the last week of January continously through the spring, summer, and fall until the last week in November. Mating and oviposition. Typically mating occurs soon after emergence from the pupa. The search for a mate is facilitated by volatile chemicals called pheromones. These chemicals are usually emitted by a virgin female and act as a sex attractant. Males detect the pheromone molecules with their antennae and fly upwind to locate the source of the chemicals, the female. The act of mating may take many hours, however, once the mating pair separates the female may begin laying fertile eggs. Pheromones are often species specific and help to reproductively isolate closely related species that occur in the same area. Females may lay eggs singly or in clusters, depending on the species. Some species, such as Orgyia antiqua, will deposit eggs on the silk surrounding the pupal skin. Other species, such as Euxoa saris, scatter eggs on the soil surface. Most species attach their eggs to the vegetation that will serve as the caterpillar host plant. For instance, Phyllodesma americana will attach a single egg to the leaf of various flowering trees that will then serve as food for the caterpillar. Egg production in species of macromoths may range from a low number of less than 100 eggs per female to a high number exceeding 1,000 eggs per female. The caterpillar is the actively feeding immature stage of moths and butterflies and is perhaps less obvious at first glance but can be abundant on certain plants at certain times of the year. Within a given environment caterpillars can be found in a variety of habitats and microhabitats. In general, they may be aquatic or terrestrial. Caterpillars can be found in fruits, roots and stems as borers or miners, in foliage as miners, on the surface of foliage as skeletonizers or chewers, in galls, or in the nests of other insects, such as ants and bees. Only the larval stage of Lepidoptera is called a caterpillar. Caterpillars initially develop in the egg and then emerge through the eggshell that they sometimes eat. The caterpillar increases in size each time it sheds its skin, a process called molting. The individual caterpillar is termed an instar between molts. Typically, a caterpillar passes through five instars as it eats and grows. In certain species a caterpillar that will become an adult female may develop through an additional instar and thus grow bigger than the male. Even into the last instar it is usually difficult to distinguish between the sexes. Most caterpillars feed and develop as solitary individuals; however, the caterpillars of a few species aggregate, some of which construct nests. For instance, the caterpillars of Lophocampa argentata aggregate on branches of Douglas-fir but do not construct nests. The caterpillars of Hyphantria cunea and Malacosoma californicum occur in large colonies living in silk nests spun across twigs and branches of trees. Caterpillar growth rates are strongly influenced by temperature and nutritional quality of host plants. Growth rates are slow at cold temperatures and up to a certain point faster at warm temperatures. Dependence of caterpillar development upon the nutritional quality of vegetation is strongly influenced by the amount of protein (nitrogen), water content, and allelochemicals. Most plants contain between 1% and 7% nitrogen by weight. Also, growth is enhanced when water content of the food is at the higher end of the normal range. Allelochemicals are plant-derived chemicals that may stimulate or deter feeding by caterpillars. Some of the better known allelochemicals are terpenes, alkaloids, phenolics, and various proteins. These chemicals may also act as poisons to the caterpillar or in certain instances the caterpillar may store poisons and in turn become toxic to potential predators. Many of the poisonous caterpillars are aposematic, that is they are brightly colored, with the colors acting as a warning signal to would be predators. For instance, the brightly colored caterpillar of the cinnabar moth, Tyria jacobaeae, is poisonous to most prospective predators due to the storage of plant derived alkaloids. The caterpillar life stage of many of the common species found in forests and woodlands of the Pacific Northwest are presented in a similar format as this book in Miller (1995). Also, the caterpillars of many species found in forests and woodlands of the eastern United States are presented in Wagner et al. (1997). Metamorphosis, the process of changing from a caterpillar into an adult occurs within the pupa. In butterflies the pupa is called a chrysalis. In moths the pupa may be covered in silk, which is called a cocoon, or the pupa may be naked but perhaps encased in rolled foliage or in the soil. When a caterpillar has attained a critical size it will change its behavior from feeding to searching for or creating a site to pupate. The pupal stage may last for 2 to 3 weeks, as in Cosmia calami, or for more than 1 year, as in Coloradia pandora. In many species the pupa overwinters. Typically, overwintering pupae are in diapause, a state of development when the eventual emergence of the adult is in an arrested condition. The adult will not mature and emerge from the pupa at the appropriate time unless the pupa is first exposed to a period of cold. Overwintering. A majority of the species of macromoths in the Pacific Northwest overwinter in the pupal stage or in the egg stage. However, some species of macromoths will overwinter in the adult stage such as the noctuid Xylem cineritia and the geometrid Triphosa haesitata. Only a few of the common species in the Pacific Northwest overwinter as a caterpillar. Some of these are the arctiids Gnophaela vermiculata, Lophocampa argentata, and Pyrrharctia isabella; the geometrid Neoalcis californiaria; and the dioptid Phryganidia californica. Natural Enemies. Lepidoptera have many natural enemies. Predators of many types devour Lepidoptera, often in great quantities. Some of the most important predators are rodents; reptiles; bats; birds; spiders; nematodes; and other insects like beetles, true bugs, and parasitoids. Also, many pathogens cause fatal diseases in Lepidoptera. Some of the most important pathogens are viruses, bacteria, protozoa, microsporidia, and fungi. Lepidoptera are equipped with defense mechanisms against such natural enemies. Physical and physiological protective features include stinging hairs in the caterpillar of Hemileuca eglanterina, camouflage, or crypsis well illustrated by the white, gray, and black tones in the forewing and hindwings of adults such as Semiothisa and Itame. Behavioral protective features include flashing bright colors or eyespots to startle predators as seen in the hindwings of the noctuid Catocala ophelia, the sphingid Paonias excaecatus, and the saturniid Antheraea polyphemus.

http://www.npwrc.usgs.gov/resource/insects/macronw/life.htm

Students will learn to determine symptoms of schizophrenia, examine the relationships among genes, neurotransmitters, and identify relevant brain structures. Describe how the behavior of a person with schizophrenia may differ from normal. Tell why schizophrenia is often considered the "cancer of mental illness." Explain why the identical twin of an individual with schizophrenia is more likely to develop the disorder than a fraternal twin or other sibling, but not guaranteed to develop it. Provide evidence that several neurotransmitters may be involved in producing schizophrenia. Elucidate the relationships among genes, neurotransmitters, and the appearance of schizophrenia. Identify brain structures that differ from normal in patients with recurring symptoms of schizophrenia. Students reveal their preconceptions about depression, then use G2C Online to learn about symptoms of the disorder, genes, and neurotransmitters associated with it, and challenges involved in diagnosis and treatment.

http://www.dnalc.org/view/1362-Schizophrenia-lesson-.html

How does it come to fossilizations? left picture shows a wall of the gorge of Bottachione, like in other gorges. We are able to recognize layers of stones. These layers are something like a geological clock. Every layer stands for a certain period of time. The top of the layers is made of young material, which has sedimented later in time (now). When new layers sedimented on the top of the older layers the pressure increases and leads to a stonification of previosly soft material. Counting the layers from the top of the earth surface therefore allows an esitmation of the age of a specific layer, similar to esitmating the age of trees by counting the age rings. When stonification of animals are found in some layers, we can roughly esitmate their age by counting the layers above the stonification. If an animal died in the past, the soft parts started to admister. The hard bones were covered in the run of the time by more and more earth material. The bones started to stonify, due to the increasing pressure and incoming Comparing the ages of all dinosaur fossil findings, it was impressive to see that no dinosaurs were found in layers younger than the border between creaceous and tertiary period. So: What caused the sudden disappearance of nearly all dinosaurs, nearly at the same time all over the world?

http://library.thinkquest.org/10981/fossilizations.html

By Andrew Bridges Vertebrates: Built on Backbones introduces students to the fascinating world of vertebrate animals. Students learn that every vertebrate has a backbone, the key feature that separates vertebrates from invertebrate animals. The book explains that a vertebrate’s internal skeleton, including its backbone, grows along with the animal. Students explore examples ranging from the family dog to mighty sharks and extinct dinosaurs as they learn about the characteristics that have allowed vertebrates to be successful in ecosystems around the world and through the ages. Students delve into the diversity of vertebrates by reading about the five main vertebrate groups—fish, amphibians, reptiles, birds, and mammals—and the distinct characteristics of each group. At the end of each two-page spread, a brief statement called The Bottom Line reinforces students’ understanding by summing up the key ideas about vertebrates covered in those pages.

http://www.sallyridestore.com/books/key-concepts-in-science/life-science/key-concepts-life-science-vertebrates.html

The role of women in society in the nineteenth century was restricted by social conventions and limited opportunities. Women were expected to marry and have children, and be financially dependent on their husbands. Schooling for girls was more limited than for boys. Women rarely had careers and most professions refused entry to women. Women were allowed to become teachers, but teaching was a low-status job, and was badly paid. During the late nineteenth century, a number of circumstances challenged women’s accepted role in society. The assumption that women should marry was complicated by a shortage of men. The limitations of schooling for women were highlighted by the Taunton Royal Commission Report on secondary education in 1868. A series of female educational pioneers emerged. Their efforts formed part of a wider movement of campaigners who sought to bring women equal rights to study, work, own property and vote.

http://www.women.qmul.ac.uk/virtual/themes/1850-1901/index.htm

English Grammar 101 English Grammar 101 Home Foreword: To the Student and Parent/Teacher Introduction: To Those Grammarians Among Us Instructions: How to Complete the Lessons Module 1: Word and Phrase Patterns Module 2: Clause Patterns Module 3: Verb Tense and Verbal Patterns Lesson 3-1: Present Tense Verbs Lesson 3-2: Present Tense Verbs (Continued) Lesson 3-3: Intensifying The Present Tense Lesson 3-4: The Present Participle as Verbal Lesson 3-5: The Present Participle as Verb and Verbal Lesson 3-6: The Present Participle as Verb and Verbal (Continued) Lesson 3-7: Remembering The Prepositional Phrase Lesson 3-8: The Present Participial Phrase Lesson 3-9: The Present Participial Phrase (Continued) Lesson 3-10: The Gerund as a Verbal Lesson 3-11: The Gerund versus The Present Participle Lesson 3-12: The Gerundial Phrase Lesson 3-13: The Gerundial Phrase versus The Participial Phrase Lesson 3-14: The Gerundial Phrase versus The Participial Phrase (Continued) Lesson 3-15: Phrases: Verb versus Gerundial versus Participial Lesson 3-16: Phrases: Verb versus Gerundial versus Participial (Continued) Quiz 3-17: Cumulative Review Lesson 3-18: The Present Infinitive as Verbal Lesson 3-19: The Infinitive Phrase as Verbal Lesson 3-20: The Infinitive Phrase as Verbal (Continued) Lesson 3-21: Identification of Verbals Lesson 3-22: Identification of Verbals (Continued) Quiz 3-23: Cumulative Review Lesson 3-24: Splitting the Infinitive Lesson 3-25: Forming the Future Tense Lesson 3-26: Past Tense Verbs Lesson 3-27: Regular and Irregular Verbs Lesson 3-28: Past Tense versus Past Participle Lesson 3-29: Present Participle versus Past Participle Lesson 3-30: Irregular Verbs Lesson 3-31: Irregular Verbs (Continued) Lesson 3-32: Irregular Verbs (Continued) Quiz 3-33: Cumulative Review Lesson 3-34: The Past Participle as Verbal Lesson 3-35: The Past Participle as Verbal (Continued) Lesson 3-36: Past Participle versus Present Participle Lesson 3-37: Past Participle versus Present Participle (Continued) Lesson 3-38: The Past Infinitive as Verbal Lesson 3-39: Past Infinitive versus Present Infinitive Lesson 3-40: Identification of Verbals Lesson 3-41: Identification of Verbals (Continued) Lesson 3-42: Identification of Verbals (Continued) Lesson 3-43: Wishing For the Future Perfect Quiz 3-44: Cumulative Review Exercise 3-45: Module 3 Self-Test Module 4: Verb Forms and Sentence Patterns Module 5: Punctuation and Capitalization Module 6: Supplement - Troublesome Words rregular Verbs (Continued) There are many irregular verbs with the same word in the past and past participle. In the paired sentences below, click to select the correct form of the verb (past or past participle) for each sentence. Since the words are the same, the first word indicates past tense and the second word indicates past participle. Please bring my coat. We have ( May I buy lunch? You ( ) lunch last time. Don't burst the bubble. He has ( ) his bubble. Catch the ball. I ( ) every ball you threw. Will we climb the hill? I have ( Do cling to your beliefs. She has ( ) to her beliefs. Please deal the cards. Have you been ( ) a good hand? Do you dive? Jim ( dived or dove dived or dove ) from the platform. Don't drag your coat. The coat has been ( ) through the mud. I dream of summer. I have never ( dreamed or dreamt dreamed or dreamt ) of winter. Please don't drown. We nearly ( ) in the storm. Fight the good fight. We have always ( ) against the odds. Will the fleas flee? They all ( ) on foot. Please don't fling your clothing. Her sweater was ( ) over her shoulder. Hang my shirt in the closet. I ( ) my coat on the chair. Do they hang people? That person was accidently ( ) to death. She did hear the noise. She has ( ) the noise before. Hold your horses! The rider ( ) the reins in his hand. You keep your promise. I have ( Will you teach the class? I ( ) last year. Copyright © 1999-2013 Cingletree Learning, LLC. All rights reserved.

http://lessons.englishgrammar101.com/EnglishGrammar101/Module3/Lesson3-32.aspx

How do the Rhizophydiales reproduce asexually? The body of a chytrid is called a thallus (plural = thalli). The thallus of a typical Rhizophydium species consists of two parts: an absorptive branching rhizoidal system that contains no nuclei; a multinucleate sporangium that ranges in shape from spherical,to oval, to pear-shaped, and to multi-lobbed (Fig. 1). (Fig. 1) Thallus consisting of a spherical sporangium and branched rhizoids (Fig. 2) Mature sporangium with two discharge pores appearing as domes, each plugged with gelatinous material Sporangium releasing a cloud of numerous zoospores through one of several discharge pores When the thallus of Rhizophydium species is fully grown, the multinucleate sporangium cleaves out unwalled, single nucleate zoospores, each bearing a single posteriorly directed flagellum. Zoospores exit the sporangium, either through one to several inoperculate openings in the sporangial wall (Fig. 2) or through a pore that opens with an operculum (a flap-like opening to the discharge pore). In some species, zoospores are released as a mass, remain quiescent at the discharge pore, and then swim away (Fig. 3). In others, zoospores are released one at a time through a discharge pore or tube. The zoospore has to use its own stored food reserves (lipids and glycogen) as it swims until it attaches to a suitable host or substrate, absorbs its flagellum, produces a wall around itself, and grows a germ tube that penetrates the substrate. This stage is called the germling (Fig. 4). The germ tube of the germling becomes the rhizoidal axis of the thallus. The rhizoidal system may be sparsely or extensively branched (Figs. 1,5). The encysted zoospore portion of the germling expands and becomes the multinucleate sporangium (Figs. 1,5). |Fig. 4 Germling consisting of encysted zoospore and tubular germ tube, which is branching||Fig 5 The zoospore cyst expands to form the sporangium and the germ tube grows and branches forming the rhizoidal system How do zoospores locate a suitable host or substrate during asexual reproduction? Zoospores of parasitic chytrids use light and chemical cues to locate hosts. Zoospores of Rhizophydium littoreum, a parasite of marine green algae, are positively phototactic toward blue light, a mechanism that might assure that zoospores swim to the photic zone where its host resides. Zoospores of both R. littoreum and B. dendrobatidis exhibit chemotaxis to specific sugars, proteins and amino acids, also a mechanism by which zoospores might detect signals to potential hosts. How do Rhizophydium species reproduce sexually? There are only a few reports of sexual reproduction among species of Rhizophydium, and they are all of algal parasites. The first report was made by Scherffel (1925) in R. granulosporum, a parasite of the green algae, Tribonema. Although several mechanisms for sexual reproduction have been described, they all involve a receptive thallus and contributing thallus or cyst. By morphological convention, the contributing thallus is considered the male gamete/ gametangium and the receptive thallus the female gamete/ gametangium. After transfer of the contents from the contributing thallus to the receptive thallus, the receptive thallus becomes a zygote and differentiates into a thick-walled resting spore containing large lipid globules in the cytoplasm. The resting spore enlarges; but the contributing thallus and receptive thallus remain distinct, often with the contributing thallus persisting as an appendage on the resting spore. After a dormancy period, the resting spore germinates directly as a sporangium or secondarily as a prosporangium, with the budding out of a sporangium from which zoospores are formed and released, continuing the life history of the chytrid. Current understanding is that meiosis occurs in the resting spore prior to germination. There are two basic schemes described for sexual reproduction among Rhizophydium species, and the differences are in the details of whether or not the receptive thallus is attached to the host substrate.

http://bama.ua.edu/~nsfpeet/Rhizophydiales%20NSF%202010/Rhizophydium%20Test%20Site/Rhizophydium%20home%20life%20history.html

To make light go backward, hold up a mirror. Light bounces off the mirror and goes back. Robert W. Boyd, a professor of optics at the University of Rochester, however, went beyond this easy, straightforward technique. Instead, in the latest example of logic-defying tricks that physicists can perform with light, Boyd and his colleagues demonstrated that an optical fiber – a glass strand that transmits pulses of light – affects light in a couple of odd ways: A pulse of light shot into the fiber departs before it enters. Within the fiber, the pulse travels backward and faster than the speed of light. Amazingly, Boyd's results do not violate any law of physics. The effect is predicted by the equations describing the propagation of waves. “This is a good example of something which is very counterintuitive that the laws of nature permit,” Boyd said. An article describing the experiment appears in the journal Science. Stopped in place In the vacuum of space, light travels at 186,171 miles per second. When it passes through a transparent material such as glass or water, it slows slightly, in effect bouncing off atoms as it moves. In 1999, physicists led by Lene Vestergaard Hau of Harvard slowed the speed of light to a leisurely 38 mph by shining it into an exotic, ultracooled material known as Bose-Einstein condensate. Two years later, Hau's group, and a second team at the Harvard-Smithsonian Center for Astrophysics, brought light to a standstill and then released it with its original properties intact. In other experiments, scientists have shown that light can at least appear to travel faster than 186,171 miles per second. Physicists hope to harness such manipulations of light to speed optical communications. One becomes three For Boyd's experiment, the scientists used an optical fiber of glass with small amounts of the metal erbium, which acts as an amplifier. A pulse of laser light was fired into the fiber. Even before the peak of the pulse entered the fiber, another pulse appeared, seemingly out of nowhere, at the far end of the fiber. This new pulse then split in two. One twin of the original pulse moved forward while the other moved backward through the fiber. The backward pulse, which traveled faster than the speed of light, and the original pulse met at the front end of the fiber, where they canceled each other. Even though one pulse momentarily became three, the experiment did not violate the law mandating conservation of energy because the amplifying effect of the erbium added a temporary surge of energy. At first glance, the experiment appears to flout the usual speed limit on the transmission of signals as the original pulse jumped to the forward-moving pulse on the other side of the fiber. However, the pulses were in a shape known as Gaussian, which is, in principle, infinite in width, though in practice not quite that wide. Thus, the outgoing pulse was actually just part of the original pulse that was reshaped by the fiber's unusual properties. “It's really kind of showing the kind of manipulation of light we can do these days,” Hau said of Boyd's experiment. Boyd said this effect might find some application in speeding optical communications, but it is, for now, mostly just an impressive trick of physics. “I find it neat,” he said. “I find it nifty.”

http://www.utsandiego.com/uniontrib/20060531/news_1c31light.html

This post originally appeared on Edutopia, a site created by the George Lucas Educational Foundation, dedicated to improving the K-12 learning process by using digital media to document, disseminate, and advocate for innovative, replicable strategies that prepare students. View Original > Just what is a game jam? It is a short event, usually only a day or two, where game developers plan, design and create a short game. Similar to a music jam session, game jams don’t involve much pre-planning and rely on immediate idea generation and improvisation. Game design companies have these jam sessions regularly, and while many of the games that happen here are digital, some are paper-based. They usually occur in one physical location to allow for immediate, organic collaboration. While there is an element of competition, most of the work is focused on collaboration towards a common goal. I have witnessed some of these game jams, and have talked with many gaming experts to learn more about it. I wanted to learn more about game jams to help me come up with ideas for how we might include an activity like this in classroom instruction. The following video will help you visualize how the process works: I see a lot of great opportunity to teach and assess 21st century skills, focus on deeper learning, and present content. Here are steps and tips to get you started: 1. Create a Flexible Space If you watch the video above, you will see the room change — literally. The room is set up in a way that allows for presentation, small group work, space for making and more. Make sure you have the space or can create a space that is fluid and can easy be transformed to meet the needs of the teams and the steps in the game jam process. 2. Provide Digital or Physical Tools As teams create and decide upon their games, they will need materials to do so. However, don’t assume that a game jam is only about creating digital games. While jammers might use a digital tool like Gamestar Mechanic to create their game, they might also use physical materials like paper, glue, magazines and scissors. Have these materials available, and provide voice and choice for students to pick what materials will best meet their needs. At a GameDesk game jam, a team developed a pizza game that was aligned to math content about fractions. This game was created with physical, not digital, materials. 3. Embrace Principles of Games Design A game jam is a great opportunity to teach fundamentals of game design, from story line and narrative to the actual mechanics. Normally, participants in game jams come to it with a lot of prior knowledge, although many game jam teams have members with very little knowledge of these mechanics. The extent to which you teach this might depend on the level of students or time constraints in the classroom. However, since the game jam is in essence a design challenge, you can align to principles of STEM or STEAM education. Make sure to give students designated time before the game jam to learn these principles. 4. Domain Analysis This is probably one of the most interesting steps of the game jam. Here, the teams investigate specific content areas (or domains) and uncover how the content is both taught and represented. For example, teams can find specific learning targets in their game content, and also note how those targets are represented visually or digitally. They research how the content is traditionally taught and also assessed. Teams craft specific learning targets from this exploration and research to ensure that, when they get to the idea phase, they can create a focused game targeting very specific learning objectives. 5. Team Building Game jams always start with team builders. In a real game jam, the team members have often never worked together before; therefore, it’s crucial to set a tone for collaboration and problem solving. It’s the same for our students. In order to set them up for success, icebreakers and other team builders need to occur. After bonding as a team and analyzing content domains, the ideation phase begins. This is where teams brainstorm and collaborate on ideas for the game itself. Relying on its collective knowledge of game design and content, the team starts to craft ideas for a game that will target a specific learning objective. These objectives are tight, and there are often not too many of them. They are created in the Domain Analysis component (step 4 above), but here the team gets to start narrowing and picking these targets as well as decide on the mechanics. 7. Deadlines and Benchmarks The game jam itself has a very specific deadline, usually a full day (eight hours) of work, including the presentations and pitches. However within the game jam day, there are further benchmarks. For example, at some point during the day, teams are no longer allowed to generate game ideas and are forced to work or “make.” This helps create the urgency for the deadline and also helps to move along the process. Consider setting specific time limits for some of the steps in the game jam. 8. Presentation, Playing and Judging The culminating event for the game jam is presenting the game product and having all participants play the game. Not only is this an important assessment, but it is also an important way to celebrate the jammers’ hard work. Judges need to have specific criteria for evaluating the games. These criteria might be different for every game jam. They might include relevance to content or curriculum, marketability, player interest, ability to collaborate, and more. As you consider a game jam for your classroom, you might focus the assessment on the content area, or simply on 21st century skills like creativity and collaboration. I know many of us have more freedom after “testing season,” so that also might be a great time to give this idea a shot. Just make sure you’re clear on the learning objectives and project outcomes that you expect from your students. Don’t forget to watch the game jammers in action in the time-lapse video above to give you a full picture of what it would look like!

http://www.andrewkmiller.com/

The value and purpose of RE RE has an important part to play as part of a broad, balanced and coherent curriculum to which all pupils are entitled. RE subject matter gives particular opportunities to promote an ethos of respect for others, to challenge stereotypes and to build understanding of other cultures and beliefs. This contributes to promoting a positive and inclusive school ethos that champions democratic values and human rights. Religious education for children and young people: - provokes challenging questions about the meaning and purpose of life, beliefs, the self, issues of right and wrong, and what it means to be human. It develops pupils’ knowledge and understanding of Christianity, other principal religions, and religious traditions that examine these questions, fostering personal reflection and spiritual development - encourages pupils to explore their own beliefs (whether they are religious or non-religious) in the light of what they learn. As they examine issues of religious belief and faith and how these impact on personal, institutional and social ethics, they express their responses, thereby building resilience to anti-democratic or extremist narratives - enables pupils to build their sense of identity and belonging which helps them flourish within their communities and as citizens in a diverse society - teaches pupils to develop respect for others including people with different faiths and beliefs, and helps to challenge prejudice - prompts pupils to consider their responsibilities to themselves and others, and to explore how they might contribute to their communities and to wider society. It encourages empathy, generosity and compassion. In summary, Religious Education is important because it helps children and young people gain wisdom in the following areas of life: - cultural, artistic, musical and literary: many great artists, composers, musicians and writers had deep religious and/or philosophical motivation and inspiration for their work. Many use religious themes and employ references to religious literature and thought in their work. How can we understand the insights they are communicating without knowledge of key religious ideas and stories? - historical and geographical, scientific and technological: what is the meaning of life? Where are we going? What is 'true'? What is ‘best’? Where do we come from? Why are people different and why do they have different tastes and preferences? What is to be gained from a diverse society? How can we understand the history and traditional cultures of Britain and other countries without a knowledge and understanding of the religious and philosophical traditions which helped form them? - moral and ethical: in the light of the many moral and ethical dilemmas we meet in life, ranging from the personal to the global, what is it to lead a good life? How do we know? Whom should we trust? How can we decide? Religious and philosophical principles and insights can help guide us when faced with moral dilemmas - personal: How can I be happy? How can I best manage my relationships? What are the skills I need to succeed in life? What emotional resources do I need to maintain a healthy lifestyle? We can get insights from religions and philosophies studied in RE and get practice in 'skills for life' such as empathy, sensitivity, humility and in thinking and communicating well - political, social and psychological: How can we best understand the relationships between people? Why do religion and belief feature in the news so much? What do religious and belief groups say about various contemporary issues? How can we best understand the religious practices and festivals celebrated by our neighbours? What motivates people? Why are our public institutions set up in the way they are? How do/should people behave when in positions of power? How do/should people react when others have power over them? Without knowledge of religions and beliefs, our gathered wisdom in all these aspects of our lives will be incomplete.

http://www.eriding.net/re/agreed_syllabus_2011/pages/aims/the-value-and-purpose-of-re.html

To teach something well and ensure that students are engaged in learning, teachers need to plan and prepare effectively. If the goal is for students is to achieve at high levels, then proper planning and preparation are expected no matter what type of teaching is going on. Project-based learning is no exception. In order to be successful, projects need to be designed with the end in mind. Without closely focused learning goals, the purpose of the project can become unclear and expectations for student learning outcomes can be miscommunicated. When designing projects, it is important to ensure that the activities planned will help your students meet the intended learning objectives. By reviewing curriculum goals, objectives and curriculum specifications, teachers make choices for establishing curricular priorities. At a very basic level, project planning involves the following steps: This simple four-step process is deceiving. Project planning is not linear; it always involves circling back to previous steps to ensure alignment. The use of Curriculum-Framing Questions and a project approach should all work together to support the learning goals and targeted curriculum specifications of the unit. Throughout the unit, there should be multiple opportunities for assessment and monitoring to measure your students’ progress. When people hear the phrase, “project-based learning”, different concepts and definitions may come to mind. These may include some of the common misconceptions below. Project-based units are long and hard to keep focused. Projects involve all kinds of ”hands-on” or “minds-on” tasks of varying complexity and length. Tasks can be as detailed and involved as a service-learning project on pollution or as simple as an in-class debate. A project will be focused as long as it is well-planned, aligned to important curriculum specifications and learning targets, and clearly states student expectations. Project-based learning means a complete change in instructional practices. Project-based learning is an instructional method in a repertoire of methods. It is not appropriate for the teaching of all skills and knowledge. It incorporates and accounts for varied teaching strategies and learning styles and is a way to build on current instruction to enrich learning experiences and make more efficient use of time. The focus of an educator has not changed. The goal remains to teach students what they need to know and need to be able to do. Project-based learning simply provides a new approach to reaching this goal. Project-based learning means a lot of work. For some teachers the shift to project-based learning may not encompass many challenges, but for others the idea may be overwhelming. If you are new to projects, it is best to start small and build upon what works well. Starting small means incorporating one or two instructional methods at a time, while building up to the complete design and implementation of a project-based unit. Starting small can mean incorporating: Little by little the benefits of project-based learning will be uncovered and the shift to projects will develop over time and lead to bigger ideas and better designs. The Assessing Projects resource provides detailed information about the benefits of student-centered assessment as well as how to use these and other assessment strategies in your classroom. See examples of teacher-created assessment plans that embed assessments throughout several different projects.

http://www.intel.my/content/www/my/en/education/k12/project-design/design/planning-projects.html

When Americans today think of slavery, many think of the antebellum South. Slavery, however, is part of New England history as well. Over the course of centuries, Americans living in the North have divorced themselves from the history of slavery in their communities. Historians and scholars are well aware of this history, but for most people it is easy to be oblivious to the legacy of slavery in the North. People in our region tend to identify with famed abolitionists like William Lloyd Garrison, or to see the historic house in the community that is rumored to be an Underground Railroad stop and as part of the virtuous, noble heritage of the North. The most arrogant among us will scorn the South for being backward enough to have adopted the dreaded institution. We are reminded of the divisions in our country’s history during each election cycle, with so-called “red states” and “blue states” that mimic the geographical patterns of slaveholding and nonslaveholding states in the nineteenth century. It is easy to forget things that are shameful about our past, but it is more important to face the facts. Thanks in part to funding from Mass Humanities, Historic Newton’s new permanent exhibition, Confronting Our Legacy: Slavery and Antislavery in the North, reveals the history of slavery in our hometown of Newton, Massachusetts. To put together this exhibition, we enlisted professional scholars as well as qualified volunteers to research the history of slavery in Newton. We learned who had been enslaved here, their names, who their enslavers were, and a little bit about what their lives were like. We also got a better sense of how Newton fit into the larger context of colonial Massachusetts and the role it played in the world of the Atlantic slave trade. On February 26, 1638, John Winthrop wrote that the ship Desire had arrived in Boston from the Caribbean with “some cotton, and tobacco, and negroes, etc.” This is the earliest record of slave importation to Massachusetts. It also describes New England’s role in the triangle trade. Slave-made products like molasses and sugar from the Caribbean were shipped to New England to make manufactured goods such as rum. Those goods could be traded in Africa for slaves or sold to Caribbean slaveholders. By 1644, Boston traders were importing slaves directly from Africa. Massachusetts, therefore, benefitted economically from the Atlantic slave trade. This three-dimensional chart illustrates the rise-and-fall of the African-American populations of Boston (shown in blue) and Newton (shown in red) over the course of the 19th century. In 1754 there were 2,711 slaves in Massachusetts, and 989 of them lived in Boston. Before 1783, when the Massachusetts Supreme Judicial Court declared slavery unconstitutional, there were at least 33 slaveholders and 50 enslaved people in Newton. Many claimed that slavery in New England was “patriarchal” and “benign” in comparison to southern slavery. Nevertheless, escapes were frequent in Massachusetts, and many slaves petitioned the General Court to end slavery or gain more rights for themselves and their children. Despite their efforts, slavery was the law of the land from 1641 to 1783. Visitors can step inside a box that provides the same amount of space that someone would have while travelling as “cargo” on a slave ship across the Middle Passage. Northern slaveholders often had only one or two slaves in their household. For whites, enslaved people provided additional labor that allowed the slaveholder to specialize in a trade or profession and to enjoy some leisure time. Slaveholding was also a status symbol in colonial America, but for the enslaved, it meant a life of hard work. In the nineteenth century, many northerners had an economic interest in slavery as well. Textile manufacturers profited from the low cost of slave-grown southern cotton, and the slave population of the southern U.S., the Caribbean, and Latin America comprised a huge market for shoes, hats, tools, and other products made in northern factories, shops, and homes. Slavery was not just a national, but a global institution, and our region took part in it in more ways than one. So what does that mean for us today? Many whose families are recent immigrants or who never took part in slaveholding may claim that this history has nothing to do with them. Some even feel attacked in discussions about slavery in the North. “This has nothing to do with me,” is a common response. But confronting the legacy of slavery is not about blame; rather, it is about acknowledging the privileges derived by many from having enslavement of Africans and African Americans as part of our national past. When we face our history, with all its sobering ugliness, and learn about the complicity of the North as well as the South in American slavery, we will be better prepared to make our world today a more socially and politically just place.

http://www.valleyadvocate.com/article_print.cfm?aid=14845

Back to Polyhedra Unit - Student Activities Students open HyperCard, make a new stack, and name it animation. Using the polygon tool they practice drawing polygons. Example: Draw a square and draw a smaller square to the right and above. Connect the vertices from one square to the other and the effect is a cube. Ask the students to draw three polyhedra using the HyperCard tools. Position them in different places on their first card. Select the three polyhedra using the lasso. Go to edit and select copy picture. Add a new card. Go to edit and select paste picture. Using the "marching ants" tool (to the left of the lasso in the toolbox), select one of the polyhedra. Go to options and select rotate. Pull on one of the dots surrounding the object and the object will rotate. The idea is to rotate each object slightly, copy/paste to a new card, rotate again and continue the process until there are about 30 cards. Return to the first card. Select the button tool (top row, middle tool in the toolbox), go to objects and select new button. Double click in the button when it appears on your card. Select script and type in the following: On MouseUpWhen you select the browse tool (top row, left in the toolbox) and click on the button you can see the animation. If you are not familiar with HyperCard the information on HyperCard Tips might be of help. Home || The Math Library || Quick Reference || Search || Help

http://mathforum.org/alejandre/workshops/hypercard.html

Reduplication is used in inflections to convey a grammatical function, such as plurality, intensification, etc., and in lexical derivation to create new words. It is often used when a speaker adopts a tone more "expressive" or figurative than ordinary speech and is also often, but not exclusively, iconic in meaning. Reduplication is found in a wide range of languages and language groups, though its level of linguistic productivity varies. Reduplication is the standard term for this phenomenon in the linguistics literature. Other terms that are occasionally used include cloning, doubling, duplication, repetition, and tautonym. Typological description Reduplication is often described phonologically in one of two different ways: either (1) as reduplicated segments (sequences of consonants/vowels) or (2) as reduplicated prosodic units (syllables or moras). In addition to phonological description, reduplication often needs to be described morphologically as a reduplication of linguistic constituents (i.e. words, stems, roots). As a result, reduplication is interesting theoretically as it involves the interface between phonology and morphology. The base is the word (or part of the word) that is to be copied. The reduplicated element is called the reduplicant, often abbreviated as RED or sometimes just R. In reduplication, the reduplicant is most often repeated only once. However, in some languages, reduplication can occur more than once, resulting in a tripled form, and not a duple as in most reduplication. Triplication is the term for this phenomenon of copying two times. Pingelapese has both reduplication and triplication. |kɔul 'to sing'||kɔukɔul 'singing'||kɔukɔukɔul 'still singing'| |mejr 'to sleep'||mejmejr 'sleeping'||mejmejmejr 'still sleeping'| Sometimes gemination (i.e. the doubling of consonants or vowels) is considered to be a form of reduplication. The term dupleme has been used (after morpheme) to refer to different types of reduplication that have the same meaning. Full and partial reduplication |[ɡin]||'ourselves'||→||[ɡinɡin]||'we (to) us'||(ɡin-ɡin)| |[jaː]||'themselves'||→||[jaːjaː]||'they (to) them'||(jaː-jaː)||(Watters 2002)| |[kʼʷə́ɬ]||'to capsize'||→||[kʼʷə́ɬkʼʷəɬ]||'likely to capsize'||(kʼʷə́ɬ-kʼʷəɬ)| |[qʷél]||'to speak'||→||[qʷélqʷel]||'talkative'||(qʷél-qʷel)||(Shaw 2004)| Partial reduplication involves a reduplication of only part of the word. For example, Marshallese forms words meaning 'to wear X' by reduplicating the last consonant-vowel-consonant (CVC) sequence of a base, i.e. base+CVC: |kagir||'belt'||→||kagirgir||'to wear a belt'||(kagir-gir)| |takin||'sock'||→||takinkin||'to wear socks'||(takin-kin)||(Moravsik 1978)| Many languages often use both full and partial reduplication, as in the Motu example below: |Base Verb||Full reduplication||Partial reduplication| |mahuta 'to sleep'||mahutamahuta 'to sleep constantly'||mamahuta 'to sleep (plural)'| Reduplicant position Initial reduplication in Agta (CV- prefix): |[ŋaŋaj]||'a long time'||→||[ŋaŋaŋaj]||'a long time (in years)'||(ŋa-ŋaŋaj)||(Healey 1960)| Final reduplication in Dakota (-CCV suffix): |[hãska]||'tall (singular)'||→||[hãskaska]||'tall (plural)'||(hãska-ska)| |[waʃte]||'good (singular)'||→||[waʃteʃte]||'good (plural)'||(waʃte-ʃte)||(Shaw 1980, Marantz 1982, Albright 2002)| Internal reduplication in Samoan (-CV- infix): |savali||'he/she walks' (singular)||→||savavali||'they walk' (plural)||(sa-va-vali)| |alofa||'he/she loves' (singular)||→||alolofa||'they love' (plural)||(a-lo-lofa)||(Moravcsik 1978, Broselow and McCarthy 1984)| |le tamaloa||'the man' (singular)||→||tamaloloa||'men' (plural)||(tama-lo-loa)| Internal reduplication is much less common than the initial and final types. Copying direction A reduplicant can copy from either the left edge of a word (left-to-right copying) or from the right edge (right-to-left copying). There is a tendency for prefixing reduplicants to copy left-to-right and for suffixing reduplicants to copy right-to-left: Final R → L copying in Sirionó: |ñimbuchao||→||ñimbuchaochao||'to come apart'||(ñimbuchao-chao)||(McCarthy and Prince 1996)| Copying from the other direction is possible although less common: Initial R → L copying in Tillamook: |[təq]||'break'||→||[qtəq]||'they break'||(q-təq)||(Reichard 1959)| Final L → R copying in Chukchi: |nute-||'ground'||→||nutenut||'ground (abs. sg.)'||(nute-nut)| |jilʔe-||'gopher'||→||jilʔejil||'gopher (abs. sg.)'||(jilʔe-jil)||(Marantz 1982)| Internal reduplication can also involve copying the beginning or end of the base. In Quileute, the first consonant of the base is copied and inserted after the first vowel of the base. Internal L → R copying in Quileute: |[tsiko]||'he put it on'||→||[tsitsko]||'he put it on (frequentative)'||(tsi-ts-ko)| |[tukoːjoʔ]||'snow'||→||[tutkoːjoʔ]||'snow here and there'||(tu-t-ko:jo’)||(Broselow and McCarthy 1984)| In Temiar, the last consonant of the root is copied and inserted before the medial consonant of the root. |[sluh]||'to shoot (perfective)'||→||[shluh]||'to shoot (continuative)'||(s-h-luh)| |[slɔɡ]||'to marry (perfective)'||→||[sɡlɔɡ]||'to marry (continuative)'||(s-ɡ-lɔɡ)||(Broselow and McCarthy 1984, Walther 2000)| A rare type of reduplication is found in Semai (an Austroasiatic language of Malaysia). "Expressive minor reduplication" is formed with an initial reduplicant that copies the first and last segment of the base: |[dŋɔh]||→||[dhdŋɔh]||'appearance of nodding constantly'||(dh-dŋɔh)| |[cruhaːw]||→||[cwcruhaːw]||'monsoon rain'||(cw-cruhaːw)||(Diffloth 1973| Reduplication and other morphological processes All of the examples above consist of only reduplication. However, reduplication often occurs with other phonological and morphological process, such as deletion, affixation of non-reduplicating material, etc. For instance, in Tz'utujil a new '-ish' adjective form is derived from other words by suffixing the reduplicated first consonant of the base followed by the segment [oχ]. This can be written succinctly as -Coχ. Below are some examples: - [kaq] 'red' → [kaqkoχ] 'reddish' (kaq-k-oχ) - [qʼan] 'yellow' → [qʼanqʼoχ] 'yellowish' (qʼan-qʼ-oχ) - [jaʔ] 'water' → [jaʔjoχ] 'watery' (jaʔ-j-oχ) (Dayley 1985) Somali has a similar suffix that is used in forming the plural of some nouns: -aC (where C is the last consonant of the base): - [toɡ] 'ditch' → [toɡaɡ] 'ditches' (toɡ-a-ɡ) - [ʕad] 'lump of meat' → [ʕadad] 'lumps of meat' (ʕad-a-d) - [wɪːl] 'boy' → [wɪːlal] 'boys' (wɪːl-a-l) (Abraham 1964) This combination of reduplication and affixation is commonly referred to as fixed-segment reduplication. - [nowiu] 'ox' → [nonnowiu] 'ox (distributive)' (no-n-nowiu) - [hódai] 'rock' → [hohhodai] 'rock (distributive)' (ho-h-hodai) - [kow] 'dig out of ground (unitative)' → [kokkow] 'dig out of ground (repetitive)' (ko-k-kow) - [ɡɨw] 'hit (unitative)' → [ɡɨɡɡɨw] 'hit (repetitive)' (ɡɨ-ɡ-ɡɨw) (Haugen forthcoming) Sometimes gemination can be analyzed as a type of reduplication. |This section requires expansion. (May 2008)| Phonological processes, environment, and reduplicant-base relations |This section requires expansion. (December 2009)| - base-reduplicant "identity" (OT terminology: BR-faithfulness) - tonal transfer/non-transfer Function and meaning In the Malayo-Polynesian family, reduplication is used to form plurals (among many other functions): - Malay rumah "house", rumah-rumah "houses". In pre-1972 Indonesian and Malay orthography, 2 was shorthand for the reduplication that forms plurals: orang "person", orang-orang or orang2 "people". This orthography has resurfaced widely in text messaging and other forms of electronic communication. Chinese also uses reduplication: 人 rén for "person", 人人 rénrén for "everybody". Japanese does it too: 時 toki "time", tokidoki 時々 "sometimes, from time to time". Both languages can use a special written iteration mark 々 to indicate reduplication, although in Chinese the iteration mark is no longer used in standard writing and is often found only in calligraphy. - spondeo, spopondi (Latin, "I vow, I vowed") - λείπω, λέλοιπα (Greek, "I leave, I left") - δέρκομαι, δέδορκα (Greek, "I see, I saw"; these Greek examples exhibit ablaut as well as reduplication) - háitan, haíháit (Gothic, "to name, I named") None of these sorts of forms survive in modern English, although they existed in its parent Germanic languages. A number of verbs in the Indo-European languages exhibit reduplication in the present stem rather than the perfect stem, often with a different vowel from that used for the perfect: Latin gigno, genui ("I beget, I begat") and Greek τίθημι, ἔθηκα, τέθηκα (I place, I placed, I have placed). Other Indo-European verbs used reduplication as a derivational process; compare Latin sto ("I stand") and sisto ("I remain"). All of these Indo-European inherited reduplicating forms are subject to reduction by other phonological laws. Contemporary spoken Finnish uses reduplicated nouns to indicate genuinity, completeness, originality and being uncomplicated as opposed to being fake, incomplete, complicated or fussy. It can be thought as compound word formation. For example, Söin viisi jäätelöä, pullapitkon ja karkkia, sekä tietysti ruokaruokaa. "I ate five choc-ices, a long loaf of coffee bread and candy, and of course food-food". Here, the "food-food" is contrasted to the "junk-food"—the principal role of food is nutrition, and "junkfood" isn't nutritious, so "food-food" is nutritious food, exclusively. One may say "En ollut eilen koulussa, koska olin kipeä. Siis kipeäkipeä" ("I wasn't at school yesterday because I was sick. Sick-sick, that is"), meaning one was actually suffering from an illness and is not making up excuses as usual. - ruoka "food", ruokaruoka "proper food", as opposed to snacks - peli "game", pelipeli "complete game",as opposed to a mod - puhelin "phone", puhelinpuhelin "phone for talking", as opposed to a pocket computer - kauas "far away", kauaskauas "unquestionably far away" - koti "home", kotikoti "home of your parents", as opposed to one's current place of residence These sorts of reduplicative forms, such as "food-food," are not merely literal translations of the Finnish but in fact have some frequency in contemporary English for emphasising, as in Finnish, an "authentic" form of a certain thing. "Food-food" is one of the most common, along with such a possibilities for "car-car" to describe a vehicle which is actually a car (small automobile) and not something else such as a truck, or "house-house," for a stand-alone house structure as opposed to an apartment, for instance. Reduplication comes after inflection in Finnish. Young adults may ask one another Menetkö kotiin vai kotiinkotiin? "Are you going home or home-home?" The reduplicated home refers to the old home that used to be their home before they moved out to their new home. In Swiss German, the verbs gah or goh "go", cho "come", la or lo "let" and aafa or aafo "begin" reduplicate when combined with other verbs. |literal translation:||she||comes||our||Christmas tree||come||adorn| |translation||She comes to adorn our Christmas tree.| |translation:||She doesn't let him sleep.| In some Salishan languages, reduplication is used to mark both diminution and plurality, one process applying to each end of the word, as in the following example from Shuswap. Note that the data was transcribed in a way that is not comparable to the IPA, but the reduplication of both initial and final portions of the root is clear: ṣōk!Emē'’n 'knife' reduplicated as ṣuk!ṣuk!Emen'’me’n 'plural small knives' (Haeberlin 1918:159). Reduplicative babbling in child language acquisition During the period 25–50 weeks after birth, all typically developing infants go through a stage of reduplicated or canonical babbling (Stark 198, Oller, 1980). Canonical babbling is characterized by repetition of identical or nearly identical consonant-vowel combinations, such as 'nanana' or 'didididi'. It appears as a progression of language development as infants experiment with their vocal apparatus and home in on the sounds used in their native language. Canonical/reduplicated babbling also appears at a time when general rhythmic behavior, such as rhythmic hand movements and rhythmic kicking, appear. Canonical babbling is distinguished from earlier syllabic and vocal play, which has less structure. The Proto-Indo-European language used partial reduplication of a consonant and e in many stative aspect verb forms. The perfect or preterite (past) tense of some Ancient Greek, Gothic, and Latin verbs preserves this reduplication: - λύω lúō 'I free' vs. λέλυκα léluka "I have freed" - hald "I hold" vs. haíhald (hĕhald) "I/he held" - currō "I run" vs. cucurrī "I ran" or "have run" Proto-Indo-European also used reduplication for imperfective aspect. Ancient Greek preserves this reduplication in the present tense of some verbs. Usually, but not always, this is reduplication of a consonant and i, and contrasts with e-reduplication in the perfect: - δίδωμι dídōmi "I give" (present) - δέδωκα dédōka "I have given" (perfect) - *σίσδω sísdō → ἵζω hízō "I set" (present) - *σέσδομαι sésdomai → ἕζομαι hézomai "I sit down" (present; from sd-, zero-grade of root in *sed-os → ἕδος hédos "seat, abode") English has several types of reduplication, ranging from informal expressive vocabulary (the first four forms below) to grammatically meaningful forms (the last two below). - Rhyming reduplication: hokey-pokey, razzle-dazzle, super-duper, boogie-woogie, teenie-weenie, walkie-talkie, wingding. Although at first glance "Abracadabra" appears to be an English rhyming reduplication it in fact is not; instead, it is derived from the Aramaic formula "Abəra kaDavəra" meaning "I would create as I spoke") - Exact reduplications (baby-talk-like): bye-bye, choo-choo, night-night, no-no, pee-pee, poo-poo. Couscous is not an English example for reduplication, since it is taken from a French word which has a Maghrebi origin. - Ablaut reduplications: bric-a-brac, chit-chat, criss-cross, ding-dong, jibber-jabber, kitty-cat, knick-knack, pitter-patter, splish-splash, zig-zag. In the ablaut reduplications, the first vowel is almost always a high vowel and the reduplicated ablaut variant of the vowel is a low vowel. - Shm-reduplication can be used with most any word; e.g. baby-shmaby, cancer-schmancer and fancy-schmancy. This process is a feature of American English from Yiddish, starting among the American Jews of New York City, then the New York dialect and then the whole country. Only the last of the above types is productive, meaning that examples of the first three are fixed forms and new forms are not easily accepted. - Comparative reduplication: In the sentence "John's apple looked redder and redder," the reduplication of the comparative indicates that the comparative is becoming more true over time, meaning roughly "John's apple looked progressively redder as time went on." In particular, this construction does not mean that John's apple is redder than some other apple, which would be a possible interpretation in the absence of reduplication, e.g. in "John's apple looked redder." With reduplication, the comparison is of the object being compared to itself over time. Comparative reduplication always combines the reduplicated comparative with "and". This construction is common in speech and is used even in formal speech settings, but it is less common in formal written texts. Although English has simple constructs with similar meanings, such as "John's apple looked ever redder," these simpler constructs are rarely used in comparison with the reduplicative form. Comparative reduplication is fully productive and clearly changes the meaning of any comparative to a temporal one, despite the absence of any time-related words in the construction. For example, the temporal meaning of "The frug seemed wuggier and wuggier" is clear: Despite not knowing what a frug is or what wugginess is, we know that the apparent wugginess of the frug was increasing over time, as indicated by the reduplication of the comparative "wuggier". - Contrastive focus reduplication: Exact reduplication can be used with contrastive focus (generally where the first noun is stressed) to indicate a literal, as opposed to figurative, example of a noun, or perhaps a sort of Platonic ideal of the noun, as in "Is that carrot cheesecake or carrot CAKE-cake?". This is similar to the Finnish use mentioned below. An extensive list of such examples is found in . More can be learned about English reduplication in Thun (1963), Cooper and Ross (1975), and Nevins and Vaux (2003). While not common in Dutch, reduplication does exist. Most, but not all (e.g., pipi, blauwblauw (laten), taaitaai (gingerbread)) reduplications in Dutch are loanwords (e.g., koeskoes, bonbon, (ik hoorde het) via via) or imitative (e.g., tamtam, tomtom). Another example is a former safe sex campaign slogan in Flanders: Eerst bla-bla, dan boem-boem (First talk, then have sex). In Dutch the verb "gaan" (to go) can be used as an auxiliary verb, which can lead to a triplication: we gaan (eens) gaan gaan (we are going to get going). The use of gaan as an auxiliary verb with itself is considered incorrect, but is commonly used in Flanders. Numerous examples of reduplication in Dutch (and other languages) are discussed by Daniëls (2000). Afrikaans regularly utilizes reduplication to emphasize the meaning of the word repeated. For example, krap means "to scratch one's self," while krap-krap-krap means "to scratch one's self vigorously." Reduplication in Afrikaans has been described extensively in the literature - see for example Botha (1988), Van Huyssteen (2004) and Van Huyssteen & Wissing (2007). Further examples of this include: "koes" (to dodge) being reduplicated in the sentence "Piet hardloop koes-koes weg" (Piet is running away while constantly dodging / cringing); "sukkel" (to struggle) becoming "sukkel-sukkel" (making slow progress; struggling on); and "kierang" (to cheat) becoming "kierang-kierang" to indicate being cheated on repeatedly . In Italian reduplication was used both to create new words or words associations (tran-tran, via via, leccalecca) and to intensify the meaning (corri!, corri! "run!, run!"). Common in Lingua Franca, particularly but not exclusively for onomatopoeic action descriptions: "Spagnoli venir...boum boum...andar; Inglis venir...boum boum bezef...andar; Francés venir...tru tru tru...chapar." ("The Spaniards came, cannonaded, and left. The English came, cannonaded heavily, and left. The French came, trumpeted on bugles, and captured it.") Common uses for reduplication in French are the creation of hypocoristics for names, whereby Louise becomes "Loulou", and Zinedine Zidane becomes Zizou; and in many infantile words, like dada, 'horse' (standard cheval), tati, 'aunt' (standard tante), or tonton, 'uncle' (standard oncle). - Romanian: mormăi, ţurţur, dârdâi, expessions talmeş-balmeş, harcea-parcea, terchea-berchea, ţac-pac, calea-valea, hodoronc-tronc, and recent slang, trendy-flendy. - Catalan: balandrim-balandram, baliga-balaga, banzim-banzam, barliqui-barloqui, barrija-barreja, bitllo-bitllo, bub-bub, bum-bum, but-but, catric-catrac, cloc-cloc, cloc-piu, corre-corrents, de nyigui-nyogui, farrigo-farrago, flist-flast, fru-fru, gara-gara, gloc-gloc, gori-gori, leri-leri, nap-buf, ning-nang, ning-ning, non-non, nyam-nyam, nyau-nyau, nyec-nyec, nyeu-nyeu, nyic-nyic, nyigo-nyigo, nyigui-nyogui, passa-passa, pengim-penjam, pif-paf, ping-pong, piu-piu, poti-poti, rau-rau, ringo-rango, rum-rum, taf-taf, tam-tam, tau-tau, tic-tac, tol·le-tol·le, tric-trac, trip-trap, tris-tras, viu-viu, xano-xano, xau-xau, xerric-xerrac, xim-xim, xino-xano, xip-xap, xiu-xiu, xup-xup, zig-zag, ziga-zaga, zim-zam, zing-zing, zub-zub, zum-zum. In colloquial Mexican Spanish it is common the use of reduplicated adverbs such as luego luego (after after) meaning "immediately", or casi casi (almost almost) which intensifies the meaning of 'almost'. Slavic languages The reduplication in the Russian language serves for various kinds of intensifying of the meaning and exists in several forms: a hyphenated or repeated word (either exact or inflected reduplication), and forms similar to shm-reduplication. Reduplication is a very common practice in Persian, to the extent that there are jokes about it. Mainly due to the mixed nature of the Persian language, most of the reduplication comes in the form of a phrase consisting of a Persian word -va- (and) and an Arabic word, like "Taghdir-Maghdir". Reduplication is particularly common in the city of Shiraz in southwestern Iran. One can further categorize the reduplicative words into "True" and "Quasi" ones. In true reduplicative words, both words are actually real words and have meaning in the language in which it is used. In quasi-reduplicative words, at least one of the words does not have a meaning. Some examples of true reduplicative words in Persian are: "Xert-o-Pert" (Odds and ends); "Čert-o-Pert" (Nonsense); "Čarand-o-Parand" (Nonsense); "Āb-o-Tāb" (much detail). Among the quasi-reduplicative words are "Zan-o-man" (wife); "Davā-Mavā" (Argument); "Talā-malā" (jewelry); and "Raxt-o-Paxt" (Items of Clothing). In general reduplication in Persian, is mainly a mockery of words with non-Persian origins. Indo-Aryan (and Dravidian) languages Typically all Indo-Aryan languages, like Hindi, Punjabi, Gujarati and Bengali use reduplication in some form or the other. It is usually used to sound casual, or in a suggestive manner. It is often used to mean etcetera. For example in Hindi, chai-shai (chai means tea, while this phrase means tea or any other supplementary drink or tea along with snacks). Quite common in casual conversations are a few more examples like shopping-wopping, khana-wana. Reduplication is also used in Dravidian languages like Telugu for the same purpose. A number of Nepalese nouns are formed by reduplication. As in other languages, the meaning is not that of a true plural, but collectives that refer to a set of the same or related objects, often in a particular situation. For example, "rangi changi"* describes an object that is extremely or vividly colorful, like a crazy mix of colors and/or patterns, perhaps dizzying to the eye. The phrase "hina mina" means "scattered," like a large collection of objects spilled (or scampering, as in small animals) in all different directions. The basic Nepalese word for food, "khana" becomes "khana sana" to refer to the broad generality of anything served at a meal. Likewise, "chiya" or tea (conventionally made with milk and sugar) becomes "chiya siya": tea and snacks (such as biscuits or cookies). *Please note, these examples of Nepalese words are spelled with a simplified Latin transliteration only, not as exact spellings. In Turkish, a word can be reduplicated while replacing the initial consonants (not being m, and possibly missing) with m. The effect is that the meaning of the original word is broadened. For example, tabak means "plate(s)", and tabak mabak then means "plates, dishes and such". This can be applied not only to nouns but to all kinds of words, as in yeşil meşil meaning "green, greenish, whatever". Although not used in formal written Turkish, it is a completely standard and fully accepted construction. Reduplication is commonly used only with 'suurensuuri' 'big of big', 'pienenpieni' 'small of small' and 'hienonhieno' 'fine of fine' but other adjectives may sometimes be duplicated as well, where a superlative is too strong an expression, somewhat similarly to Slavic languages. The structure may be written also separately as 'genitive' 'nominative', which may create confusion on occasion (f.e. 'suurensuuri jalka' 'big of big foot' vs. 'suuren suuri jalka' 'big foot of a big one') Reduplication is usually rhyming. It can add emphasis: 'pici' (tiny) -> ici-pici (very tiny) and it can modify meaning: 'néha-néha' ('seldom-seldom': seldom but repeatedly), 'erre-arra' ('this way-that way', meaning movement without a definite direction), 'ezt-azt' ('this-that', meaning 'all sort of things'), Reduplication often evokes a sense of playfulness and it's quite common when talking to small children. Bantu languages - Swahili piga 'to strike'; pigapiga 'to strike repeatedly' - Ganda okukuba (oku-kuba) 'to strike'; okukubaakuba (oku-kuba-kuba) 'to strike repeatedly, to batter' - Chewa tambalalá 'to stretch one's legs'; tambalalá-tambalalá to stretch one's legs repeatedly' Popular names that have reduplication include Semitic languages frequently reduplicate consonants, though often not the vowels that appear next to the consonants in some verb form. This can take the shape of reduplicating the antepenultimate consonant (usually the second of three), the last of two consonants, or the last two consonants. In the Hebrew, reduplication is used in nouns and adjectives. For stress, as in גבר גבר (Gever Gever) where the noun גבר 'man' - is duplicated to mean a manly man, a man among man. Or as in לאט לאט (le-aht le-aht) where the adverb לאט 'slowly' - is duplicated to mean very slowly. Meaning every, as in יום יום (yom yom) where the noun יום 'day' is duplicated to every day, day in day out, day by day. Some nouns and adjectives can also be made into diminutives by reduplication of the last two consonants (biconsonantal reduplication), e.g. - כלב (Kelev) = Dog - כלבלב (Klavlav) = Puppy - חתול (Chatul) = Cat - חתלתול (Chataltul) = Kitten - לבן (Lavan) = White - לבנבן (Levanban) = Whitish - קטן (Katan) = Small - קטנטן (Ktantan) = Tiny Reduplication in Hebrew is also productive for the creation of verbs, by reduplicating the root or part of it e.g.: dal (דל) 'poor,spare' > dilel (דלל) 'to dilute' but also dildel (דלדל) 'to impoverish, to weaken'; nad (נד) 'to move, to nod' > nadad (נדד) 'to wander' but also nidned (נדנד) 'to swing, to nag'. In Amharic, verb roots can be reduplicated three different ways. These can result in verbs, nouns, or adjectives (which are often derived from verbs). From the root sbr 'break', antepenultimate reduplication produces täsäbabbärä 'it was shattered' and biconsonantal reduplication produces täsbäräbbärä 'it was shattered repeatedly' and səbərbari 'a shard, a shattered piece'. From the root kHb 'pile stones into a wall', since the second radical is not fully specified, what some call "hollow", the antepenultimate reduplication process reduplicates the k, which is by some criteria antepenultimate, and produces akakabä 'pile stones repeatedly'. In Burmese, reduplication is used in verbs and adjectives to form adverbs. Many Burmese words, especially adjectives such as လှပ ('beautiful' [l̥a̰pa̰]), which consist of two syllables (when reduplicated, each syllable is reduplicated separately), when reduplicated (လှပ → လှလှပပ 'beautifully' [l̥a̰l̥a̰ pa̰pa̰]) become adverbs. This is also true of many Burmese verbs, which become adverbs when reduplicated. Some nouns are also reduplicated to indicate plurality. For instance, ပြည်, means "country," but when reduplicated to အပြည်ပြည်, it means "many countries" (as in အပြည်ပြည်ဆိုင်ရာ, "international"). Another example is အမျိုး, which means "kinds," but the reduplicated form အမျိုးမျိုး means "multiple kinds." A few measure words can also be reduplicated to indicate "one or the other": - ယောက် (measure word for people) → တစ်ယောက်ယောက် (someone) - ခု (measure word for things) → တစ်ခုခု (something) Adjective reduplication is common in Standard Chinese, typically denoting emphasis, less acute degree of the quality described, or an attempt at more indirect speech: xiǎoxiǎo de 小小的 (small), chòuchòu de 臭臭的 (smelly) (this can also reflect a "cute", juvenile or informal register). In the case of adjectives composed of two characters (morphemes), generally each of the two characters is reduplicated separately: piàoliang 漂亮 (beautiful) reduplicates as piàopiàoliangliang 漂漂亮亮. Verb reduplication is also common in Standard Chinese, conveying the meaning of informal and temporary character of the action. It is often used in imperative expressions, in which it lessens the degree of imperativity: zuòzuò 坐坐 (sit (for a while)), děngděng 等等 (wait (for a while)). Compound verbs are reduplicated as a whole word: xiūxixiūxi 休息休息 (rest (for a while)). This can be analyzed as an instance of omission of "一" (originally, e.g., "坐一坐" or "等一等" ) or "一下" (originally, e.g., "坐一下"). Noun reduplication, though nearly absent in Standard Chinese, is found in the southwestern dialect of Mandarin. For instance, in Sichuan Mandarin, bāobāo 包包 (handbag) is used whereas Beijing use bāor 包儿 (one exception is the colloquial use of bāobāo 包包 by non-Sichuan Mandarin speakers to reflect a perceived fancy or attractive purse). However, there are few nouns that can be reduplicated in Standard Chinese, and reduplication denotes generalisation and uniformity: rén 人 (human being) and rénrén 人人 (everybody (in general, in common)), jiājiāhùhù 家家户户 (every household (uniformly)) - in the latter jiā and hù additionally duplicate the meaning of household, which is a common way of creating compound words in Chinese. A small number of native Japanese nouns have collective forms produced by reduplication (possibly with rendaku), such as 人々 hitobito "people" (h → b is rendaku) – these are written with the iteration mark "々" to indicate duplication. This formation is not productive and is limited to a small set of nouns. Similarly to Standard Chinese, the meaning is not that of a true plural, but collectives that refer to a large, given set of the same object; for example, the formal English equivalent of 人々 would be "people" (collective), rather than "persons" (plural individuals). Japanese also contains a large number of mimetic words formed by reduplication of a syllable. These words include not only onomatopoeia, but also words intended to invoke non-auditory senses or psychological states. By one count, approximately 43% of Japanese mimetic words are formed by full reduplication, and many others are formed by partial reduplication, as in がささ〜 ga-sa-sa- (rustling) – compare English "a-ha-ha-ha". Words called từ láy are found abundantly in Vietnamese. They are formed by repeating a part of a word to form new words, altering the meaning of the original word. Its effect is to sometimes either increase or decrease the intensity of the adjective, and is often used as a literary device (like alliteration) in poetry and other compositions, as well as in everyday speech. Examples of reduplication increasing intensity: - đau → đau điếng: hurt → hurt horribly - mạnh → mạnh mẽ: strong → very strong - rực → rực rỡ: flaring → blazing Examples of reduplication decreasing intensity: - nhẹ → nhè nhẹ: soft → soft (less) - xinh → xinh xinh: pretty → cute - đỏ → đo đỏ: red → somewhat red - xanh → xanh xanh: blue/green → somewhat blue/green Examples of blunt sounds or physical conditions: - loảng xoảng — sound of glass breaking to pieces or metallic objects falling to the ground - hớt hơ hớt hải- (also hớt ha hớt hải) — hard gasps -> in extreme hurry, in panic, panic-stricken - lục đục — the sound of hard, blunt (and likely wooden) objects hitting against each other -> disagreements and conflicts inside a group or an organisation Khmer uses reduplication for several purposes, including emphasis and pluralization. Reduplication in Khmer, like many Mon–Khmer languages, can express complex thoughts. Khmer also uses a form of reduplication known as "synonym compounding", in which two phonologically distinct words with similar or identical meanings are combined, either to form the same term or to form a new term altogether. The wide use of reduplication is certainly one of the most prominent grammatical features of Indonesian and Malay (as well as of other South-East Asian and Austronesian languages). Malay and Indonesian In Malay and Indonesian, reduplication is a very productive process. It is used for expression of various grammatical functions (such as verbal aspect) and it is part in a number of complex morphological models. Simple reduplication of nouns and pronouns can express at least 3 meanings: - Diversity or non-exhaustive plurality: - Burung-burung itu juga diekspor ke luar negeri = "All those birds are also exported out of the country". - Conceptual similarity: - langit-langit = "ceiling; palate; etc." < langit = "sky"; - jari-jari = "spoke; bar; radius; etc." < jari = "finger" etc. - Pragmatic accentuation: - Saya bukan anak-anak lagi! "I am not a child anymore!" (anak = "child") Reduplication of an adjective can express different things: - Adverbialisation: Jangan bicara keras-keras! = "Don't speak loudly!" (keras = hard) - Plurality of the corresponding noun: Rumah di sini besar-besar = "The houses here are big" (besar = "big"). Reduplication of a verb can express various things: - Simple reduplication: - Pragmatic accentuation: Kenapa orang tidak datang-datang? = "Why aren't people coming?" - Reduplication with me- prefixation, depending on the position of the prefix me-: - Repetition or continuation of the action: Orang itu memukul-mukul anaknya: "That man continuously beat his child"; - Reciprocity: Kedua orang itu pukul-memukul = "Those two men would beat each other". Notice that in the first case, the nasalisation of the initial consonant (whereby /p/ becomes /m/) is repeated, while in the second case, it only applies in the repeated word. Reduplication can convey a simple plural meaning, for instance wahine "woman", waahine "women", tangata "person", taangata "people". Biggs calls this "infixed reduplication". It occurs a small subset of people words in most Polynesian languages. Reduplication can convey emphasis or repetition, for example mate "die", matemate "die in numbers"; and de-emphasis, for example wera "hot" and werawera "warm". Reduplication can also extend the meaning of a word; for instance paki "pat" becomes papaki "slap or clap once" and pakipaki "applaud"; kimo "blink" becomes kikimo "close eyes firmly". In Japanese imperative oit'oi'te (leave behind) of the compound verb oitoku is pseudo-reduplication. It appears to be 'oit' repeated, especially when spoken quickly, but the root is 'oite'(leave)+'oku'(to place something). Therefore the 'oit' sound is repeated twice, but its by chance placement and not repetition (different meanings). Australian Aboriginal languages Reduplication is common in many Australian place names due to their Aboriginal origins. Examples: Turramurra, Parramatta, Wooloomooloo. In the language of the Wiradjuri people of south eastern Australian, plurals are formed by doubling a word, hence 'Wagga' meaning crow becomes Wagga Wagga meaning 'place of many crows'. This occurs in other place names deriving from the Wiradjuri language including Gumly Gumly, Grong Grong and Book Book. See also - Language acquisition - Syntactic doubling - For an example of a language with many types of reduplication see: St'at'imcets language#Reduplication. - Word word - List of people with reduplicated names - Pratt, George (1984) . A Grammar and Dictionary of the Samoan Language, with English and Samoan vocabulary (3rd and revised ed.). Papakura, New Zealand: R. McMillan. ISBN 0-908712-09-X. Retrieved 8 June 2010. - The Malay Spelling Reform, Asmah Haji Omar, (Journal of the Simplified Spelling Society, 1989-2 pp.9-13 later designated J11) - Jila Ghomeshi, Ray Jackendoff, Nicole Rosen, and Kevin Russell (2004). "Contrastive focus reduplication in English (the Salad-Salad paper)". Natural Language & Linguistic Theory 22 (2): 307–357. doi:10.1023/B:NALA.0000015789.98638.f9. - A Glossary of Lingua Franca, 5th ed. - Peter Unseth. 2003. Surveying bi-consonantal reduplication in Semitic. In Selected Comparative-Historical Afrasian Linguistic Studies in Memory of Igor M. Diakonoff, ed. by M. Lionel Bender, 257-273. Munich: Lincom Europa. - p. 1029. Wolf Leslau. 1995. Reference Grammar of Amharic. Wiesbaden: Harrassowitz. - Peter Unseth. 2002. Biconsonantal reduplication in Amharic. Doctoral dissertation, University of Texas at Arlington. - p. 1035. Wolf Leslau. 1995. Reference Grammar of Amharic. Wiesbaden: Harrassowitz. - Tamamura, Fumio. 1979. Nihongo to chuugokugo ni okeru onshoochoogo [Sound-symbolic words in Japanese and Chinese]. Ootani Joshidai Kokubun 9:208-216. - Tamamura, Fumio. 1989. Gokei [Word forms]. In Kooza nihongo to nihongo kyooiku 6, ed. Fumio Tamamura, 23-51. Tokyo: Meiji Shoin. - Reduplicants and Prefixes in Japanese Onomatopoeia, Akio Nasu - Yury A. Lande, "Nominal reduplication in Indonesian challenging the theory of grammatical change", International Symposium on Malay/Indonesian Linguistics, Nijmegen, The Netherlands, 27–29 June 2003. - Biggs, Bruce, 1998. Let's learn Maori: a guide to the study of the Maori language. Auckland: Auckland University Press, p137. - Abraham, Roy. (1964). Somali-English dictionary. London, England: University of London Press. - Albright, Adam. (2002). A restricted model of UR discovery: Evidence from Lakhota. (Draft version). - Alderete, John; Benua, Laura; Gnanadesikan, Amalia E.; Beckman, Jill N.; McCarthy, John J.; and Urbanczyk, Suzanne. (1999). Reduplication with fixed segmentism. Linguistic Inquiry, 30, 327-364. (Online version ROA 226-1097). - Botha, Rudi P. (1988). Form and meaning in word formation : a study of Afrikaans reduplication. Cambridge: Cambridge University Press. - Broselow, Ellen; and McCarthy, John J. (1984). A theory of internal reduplication. The linguistic review, 3, 25-88. - Cooper, William E.; and Ross, "Háj" John R. (1975). World order. In R. E. Grossman, L. J. San, and T. J. Vance (Eds.), Papers from the parasession on functionalism (pp. 63–111). Chicago, IL: Chicago Linguistic Society. - Dayley, Jon P. (1985). Tzutujil grammar. Berkeley, CA: University of California Press. - Diffloth, Gérald. (1973). Expressives in Semai. In P. N. Jenner, L. C. Thompson, and S. Starsota (Eds.), Austroasiatic studies part I (pp. 249–264). University Press of Hawaii. - Fabricius, Anne H. (2006). A comparative survey of reduplication in Australian languages. LINCOM Studies in Australian Languages (No. 03). Lincom. ISBN 3-89586-531-1. - Haeberlin, Herman. (1918). “Types of Reduplication in Salish Dialects.” International Journal of American Linguistics 1: 154-174. - Haugen, Jason D. (forthcoming). Reduplicative allomorphy and language prehistory in Uto-Aztecan. (Paper presented at Graz Reduplication Conference 2002, November 3–6). - Harlow, Ray. (2007) Māori: a linguistic introduction Cambridge University Press. ISBN 978-0-521-80861-3. 127-129 - Healey, Phyllis M. (1960). An Agta grammar. Manila: The Institute of National Language and The Summer Institute of Linguistics. - Hurch, Bernhard (Ed.). (2005). Studies on reduplication. Empirical approaches to language typology (No. 28). Mouton de Gruyter. ISBN 3-11-018119-3. - Inkelas, Sharon; & Zoll, Cheryl. (2005). Reduplication: Doubling in morphology. Cambridge studies in linguistics (No. 106). Cambridge University Press. ISBN 0-521-80649-6. - Key, Harold. (1965). Some semantic functions of reduplication in various languages. Anthropological Linguistics, 7(3), 88-101. - Marantz, Alec. (1982). Re reduplication. Linguistic Inquiry 13: 435-482. - McCarthy, John J. and Alan S. Prince. (1986 ). Prosodic morphology 1986. Technical report #32. Rutgers University Center for Cognitive Science. (Unpublished revised version of the 1986 paper available online on McCarthy's website: http://ruccs.rutgers.edu/pub/papers/pm86all.pdf). - McCarthy, John J.; and Prince, Alan S. (1995). Faithfulness and reduplicative identity. In J. Beckman, S. Urbanczyk, and L. W. Dickey (Eds.), University of Massachusetts occasional papers in linguistics 18: Papers in optimality theory (pp. 249–384). Amherst, MA: Graduate Linguistics Students Association. (Available online on the Rutgers Optimality Archive website: http://roa.rutgers.edu/view.php3?id=568). - McCarthy, John J.; and Prince, Alan S. (1999). Faithfulness and identity in prosodic morphology. In R. Kager, H. van der Hulst, and W. Zonneveld (Eds.), The prosody morphology interface (pp. 218–309). Cambridge: Cambridge University Press. (Available online on the Rutgers Optimality Archive website: http://roa.rutgers.edu/view.php3?id=562). - Moravcsik, Edith. (1978). Reduplicative constructions. In J. H. Greenberg (Ed.), Universals of human language: Word structure (Vol. 3, pp. 297–334). Stanford, CA: Stanford University Press. - Nevins, Andrew; and Vaux, Bert. (2003). Metalinguistic, shmetalinguistic: The phonology of shm-reduplication. (Presented at the Chicago Linguistics Society, April 2003). (Online version: http://ling.auf.net/lingbuzz/@qclBWVDkyQupkDAI/yuTibEgY?78). - Oller, D. Kimbrough. 1980. The emergence of the sounds of speech in infancy, in Child Phonology Vol. I, edited by G. H. Yeni-Komshian, J. F. Kavanaugh, and C. A. Ferguson. Academic Press, New York. pp. 93–112. - Raimy, Eric. (2000). Remarks on backcopying. Linguistic Inquiry 31:541-552. - Rehg, Kenneth L. (1981). Ponapean reference grammar. Honolulu: The University Press of Hawaii. - Reichard, Gladys A. (1959). A comparison of five Salish languages. International Journal of American Linguistics, 25, 239-253. - Shaw, Patricia A. (1980). Theoretical Issues in Dakota Phonology and Morphology. Garland Publ: New York. pp. ix + 396. - Shaw, Patricia A. (2004). Reduplicant order and identity: Never trust a Salish CVC either?. In D. Gerdts and L. Matthewson (Eds.), Studies in Salish linguistics in honor of M. Dale Kinkade. University of Montana Occasional Papers in Linguistics (Vol. 17). Missoula, MT: University of Montana. - Stark, Rachel E. (1980). Features of infant sounds: The emergence of cooing. Journal of Child Language Vol 5(3) Oct 1978, 379-390. - Thun, Nils. (1963). Reduplicative words in English: A study of formations of the types tick-tock, hurly-burly, and shilly-shally. Uppsala. - Van Huyssteen, Gerhard B. (2004). Motivating the composition of Afrikaans reduplications: a cognitive grammar analysis. In: Radden, G & Panther, K-U. (eds.). Studies in Linguistic Motivation. ISBN 3-11-018245-9. Berlin: Mouton de Gruyter. pp. 269–292. - Van Huyssteen, Gerhard B and Wissing, Daan P. (2007). Datagebaseerde Aspekte van Afrikaanse Reduplikasies. [Data-based Aspects of Afrikaans Reduplications]. Southern African Linguistics and Applied Language Studies. 25(3): 419–439. - Watters, David E. (2002). A grammar of Kham. Cambridge grammatical descriptions. Cambridge: Cambridge University Press. ISBN 0-521-81245-3. - Wilbur, Ronnie B. (1973). The phonology of reduplication. Doctoral dissertation, University of Illinois. (Also published by Indiana University Linguistics Club in 1973, republished 1997.) |Look up reduplication in Wiktionary, the free dictionary.| - Reduplication (Lexicon of Linguistics) - What is reduplication? (SIL) - Echo-Word Reduplication Lexicon - Exhaustive list of reduplications in English - List of contrastive focus reduplications in English - graz database on reduplication (gdr) Institute of Linguistics, University of Graz - La réduplication à m dans l’arabe parlé à Mardin

http://en.wikipedia.org/wiki/Reduplication

The males and females of many bird species are difficult to distinguish by their appearance (peacocks are a notable exception). There are many situations where it is useful to know the sex of birds including captive breeding programmes, behavioural studies and even species delimitation in extinct taxa. DNA sexing provides a simple and quick way to determine which birds are females and which are males. We have been using this technique for some of our bird research projects, including our study of the prion wreck of 2011. For our prion study we want to determine whether there is a gender bias in the birds that were wrecked. So how does DNA sexing work for birds? By way of background, birds have a different chromosome system to us for determining their sex. In mammals, including us, males have an X and a Y chromosome and females have two X chromosomes. In contrast, birds have a ZW sex-determination system whereby males have two Z chromosomes and females both Z and W chromosomes. To genetically sex a bird, DNA is first obtained from a blood, feather or tissue sample. We used tongue samples for the prions. From these DNA samples we made lots of copies of the CHD region, a gene that occurs on both the Z and W chromosomes. Our processing of these gene copies produces a single DNA band for males (because they only have one type of chromosome) and two bands for females (representing the different CHD copies from the Z and W chromosomes). DNA sexing is also possible for humans, albeit using a modified method suited to our X/Y chromosome system, and is routinely used in forensics. A recent example is the detection of female DNA on the bombs used in the Boston marathon bombing.

http://blog.tepapa.govt.nz/author/shepherdlara/

Gold in California Gold in California became highly concentrated there as the result of global forces operating over hundreds of millions of years. Volcanoes, tectonic plates and erosion all combined to concentrate billions of dollars worth of gold in the mountains of California. During the California Gold Rush, gold-seekers known as "Forty-Niners" retrieved this gold, at first using simple techniques, and then developing more sophisticated techniques, which spread around the world. Scientists believe that over a span of at least 400 million years, gold that had been widely dispersed in the Earth’s crust became more concentrated by geologic actions into the gold-bearing regions of California. Only gold that is concentrated can be economically recovered. Some 400 million years ago, California lay at the bottom of a large sea; underwater volcanoes deposited lava and minerals (including gold) onto the sea floor; sometimes enough that islands were created. Between 400 million and 200 million years ago, geologic movement forced the sea floor and these volcanic islands and deposits eastwards, colliding with the North American continent, which was moving westwards. Beginning about 200 million years ago, tectonic pressure forced the sea floor beneath the American continental mass. As it sank, or subducted, beneath today's California, the sea floor heated and melted into very large molten masses (magma). Being lighter and hotter than the ancient continental crust above it, this magma forced its way upward, cooling as it rose to become the granite rock found throughout the Sierra Nevada and other mountains in California today — such as the sheer rock walls and domes of Yosemite Valley. As the hot magma cooled, solidified, and came in contact with water, minerals with similar melting temperatures tended to concentrate themselves together. As it solidified, gold became concentrated within the magma, and during this cooling process, veins of gold formed within fields of quartz because of the similar melting temperatures of both. As the Sierra Nevada and other mountains in California were forced upwards by the actions of tectonic plates, the solidified minerals and rocks were raised to the surface and exposed to rain, ice and snow. The surrounding rock then eroded and crumbled, and the exposed gold and other materials were carried downstream by water. Because gold is denser than almost all other minerals, this process further concentrated the gold as it sank, and pockets of gold gathered in quiet gravel beds along the sides of old rivers and streams. The California mountains rose and shifted several times within the last fifty million years, and each time, old streambeds moved and were dried out, leaving the deposits of gold resting within the ancient gravel beds where the gold had been collecting. Newer rivers and streams then developed, and some of these cut through the old channels, carrying the gold into still larger concentrations. The Forty-Niners of the California Gold Rush first focused their efforts on these deposits of gold, which had been gathered in the gravel beds by hundreds of millions of years of geologic action. Gold recovery The early Forty-Niners panned for gold in California’s rivers and streams, or used "cradles" and "rockers" or "long-toms," forms of placer mining. Modern estimates by the U.S. Geological Survey are that some 12 million ounces (373 t) of gold were removed in the first five years of the Gold Rush (worth approximately US$7.2 billion at November 2006 prices). By 1853, the first hydraulic mining was used. In hydraulic mining, (which was invented in California) a powerful stream of water is directed at gold-bearing gravel beds; the gravel and gold then pass over sluices, with the gold settling to the bottom. By the mid-1880s, it is estimated that 11 million ounces (342 t) of gold (worth approximately US$6.6 billion at November 2006 prices) had been recovered via "hydraulicking." The final stage to recover loose gold was to prospect for gold in the flat rivers of California’s Central Valley and other gold-bearing areas of California (such as Scott Valley in Siskiyou County). By the late 1890s, dredging technology (which was also invented in California) had become economical, and it is estimated that more than 20 million ounces (622 t) were recovered by dredging (worth approximately US$12 billion at November 2006 prices). Gold-seekers also engaged in "hard-rock" mining, that is, extracting the gold directly from the rock that contained it (typically quartz) Once the gold-bearing rocks were brought to the surface, the rocks were crushed, and the gold was separated out (using moving water), or leached out, typically by using arsenic or mercury. Eventually, hard-rock mining wound up being the single largest source of gold produced in the Gold Country. Geological after-effects There were decades of minor earthquakes, more than at any other time in the historical record for Northern California, before the 1906 San Francisco earthquake. Previously interpreted as precursory activity to the 1906 earthquake, they have been found to have a strong seasonal pattern and due to large seasonal sediment loads in coastal bays that overlie faults as a result of mining of gold inland. - Hill, Mary (1999). Gold: the California story. Berkeley and Los Angeles: University of California Press. p. 167. - Hill, Mary (1999), p. 168. - Hill, Mary (1999), pp. 168-69. - Brands, H.W. (2003). The age of gold: the California Gold Rush and the new American dream. New York: Doubleday., pp. 195-196. - Hill, Mary (1999), pp. 149-58. Similar forces created the granite domes and spires of Castle Crags in Shasta County. - Hill, Mary (1999), pp. 174-78. - Hill, Mary (1999), pp. 169-173. - Hill, Mary (1999), pp. 94-100. - Hill, Mary (1999), pp. 105-110. - Images and detailed description of placer mining tools and techniques; image of a long tom - Brands, H.W. (2002), pp. 198-200. - Bancroft, Hubert Howe (1884-1890) History of California, vols. 18-24, The works of Hubert Howe Bancroft, complete text online, pp. 87-88. - Mining History and Geology of the Mother Lode (accessed Oct. 16, 2006) - Starr, Kevin (2005). California: a history. New York: The Modern Library., p. 89. - * Rawls, James J. and Orsi, Richard J. (eds.) (1999). A golden state: mining and economic development in Gold Rush California (California History Sesquicentennial Series, 2). Berkeley and Los Angeles: Univ. of California Press. pp. 1991 - Rawls, James J. and Orsi, Richard (eds.) (1999), pp. 36-39 - Rawls, James J. and Orsi, Richard (eds.) (1999), pp. 39-43 - Seasonal Seismicity of Northern California Before the Great 1906 Earthquake, (Journal) Pure and Applied Geophysics, ISSN 0033-4553 (Print) 1420-9136 (Online), volume 159, Numbers 1-3 / January, 2002, Pages 7-62. - Bancroft, Hubert Howe (1884–1890) History of California, vols. 18-24, The works of Hubert Howe Bancroft, complete text online - Brands, H.W. (2003). The age of gold: the California Gold Rush and the new American dream. New York City: Doubleday. ISBN 0-385-72088-2. - Hill, Mary (1999). Gold: the California story. Berkeley and Los Angeles: University of California Press. ISBN 0-520-21547-8. - Rawls, James J. and Orsi, Richard J. (eds.) (1999). A golden state: mining and economic development in Gold Rush California (California History Sesquicentennial Series, 2). Berkeley and Los Angeles: University of California Press. ISBN 0-520-21771-3. - Starr, Kevin (2005). California: a history. New York: The Modern Library. ISBN 0-679-64240-4.

http://en.wikipedia.org/wiki/Gold_in_California

On this day in 1961, President John F. Kennedy issued a statement extending his "sincere wishes" and those of the American people to Soviet Premier Nikita Khrushchev and the people of the Soviet Union for a peaceful and prosperous New Year. It was the height of the Cold War and the United States and Soviet Union were locked in a nuclear arms race. Citing 1961 as a "troubled one" between the two superpowers, Kennedy said that it was his "earnest hope" that 1962 would see improved relations between the two countries. Kennedy then told Khrushchev he believed the responsibility to achieve world peace rested on the two men's shoulders. Kennedy's message came in response to a December 29 message from Khrushchev that carried his hope that 1962 would be a "threshold" year for taking "efficient steps in the cause of liquidation of centers of military danger." Khrushchev was likely referring to tensions over the ongoing division of the city of Berlin into democratic and communist sectors. In August 1961, it was Khrushchev's government that approved East Germany's decision to construct a physical barrier, the Berlin Wall, between the two sectors to stop communist-ruled East Germans from defecting to the West. Although Kennedy and Khrushchev both pledged cooperation as 1961 came to a close, the two went on to play a dangerous game of chicken over Soviet missile sites in Cuba in October 1962, leading the world to the very brink of nuclear war during the Cuban Missile Crisis.

http://www.history.com/this-day-in-history/kennedy-and-khrushchev-exchange-holiday-greetings?catId=9

U.S. Supreme Court Justice Thurgood Marshall was born July 2, 1908 in Baltimore, Maryland. He received his B.A. from Lincoln University as a pre-dental student. Later, he changed his major to law and graduated magna cum laude from Howard University's Law School in 1933. As an attorney in private practice, Marshall defended people who were too poor to afford legal fees. In 1936, he began his career with the NAACP, and eventually became director of its Legal Defense Fund. In 1946, Marshall was presented the NAACP Spingarn Medal for his services as a lawyer appearing before the Supreme Court. However, in 1954 as part of an imposing team of lawyers, Marshall gave the closing arguments to probably the most monumental case of his career, Brown vs. Board of Education. Marshall spoke eloquently, and asked: "Why of all the multitudinous groups of people in this country [do] you have to single out Negroes and give them this separate treatment?" The Supreme Court ruled "separate but equal" to be unconstitutional. On June 13, 1967, at age 59, Thurgood Marshall became the ninety-sixth man, and the first Black to be appointed to the highest court in the land, the U.S. Supreme Court.

http://culturebus.com/site/?page=profile&url_id=173&n=Marshall,_Thurgood

Coagulation or clotting stops bleeding. Coagulation is a complex process involving several steps. If a blood vessel is damaged, thrombocytes (platelets) line the wall of the blood vessel. The platelets clump together. This process is called aggregation. Clotting factors, which are particular proteins formed in the liver, also travel to the wounded blood vessel. A complex chain reaction including these factors gathers more platelets and repairs the wound. The walls of the wound close together and connective tissue cells help build new tissue.

http://www.informedhealthonline.org/sidgioqoktsq5l79e2q4f9f5pgv0e416393ik/dictionary.57.en.html?bab[mode]=abc&bab[entry_id]=240&bab[search_offset]=20

Detecting Massive Neutrinos; August 1999; Scientific American Magazine; by Kearns, Kajita, Totsuka; 8 Page(s) One man's trash is another man's treasure. For a physicist, the trash is "background"-some unwanted reaction, probably from a mundane and well-understood process. The treasure is "signal"-a reaction that we hope reveals new knowledge about the way the universe works. Case in point: over the past two decades, several groups have been hunting for the radioactive decay of the proton, an exceedingly rare signal (if it occurs at all) buried in a background of reactions caused by elusive particles called neutrinos. The proton, one of the main constituents of atoms, seems to be immortal. Its decay would be a strong indication of processes described by Grand Unified Theories that many believe lie beyond the extremely successful Standard Model of particle physics. Huge proton-decay detectors were placed deep underground, in mines or tunnels around the world, to escape the constant rain of particles called cosmic rays. But no matter how deep they went, these devices were still exposed to penetrating neutrinos produced by the cosmic rays. The first generation of proton-decay detectors, operating from 1980 to 1995, saw no signal, no signs of proton decay-but along the way the researchers found that the supposedly mundane neutrino background was not so easy to understand. One such experiment, Kamiokande, was located in Kamioka, Japan, a mining town about 250 kilometers (155 miles) from Tokyo (as the neutrino flies). The name stood for "Kamioka Nucleon Decay Experiment." Scientists there and at the IMB experiment, located in a salt mine near Cleveland, Ohio, used sensitive detectors to peer into ultrapure water, waiting for the telltale flash of a proton decaying.

http://www.sciamdigital.com/index.cfm?fa=Products.ViewIssuePreview&ARTICLEID_CHAR=D3F26459-6BF8-4427-BBAE-FCD1496E544

Skip to main content More Search Options A member of our team will call you back within one business day. The middle ear is an air-filled chamber that lies behind the eardrum. Pressure in the middle ear changes to match air pressure outside of the eardrum. When inside and outside pressures are balanced, the eardrum is flexible and normal hearing is more likely. Problems occur when air pressure in the middle ear drops. This is usually due to a block in the eustachian (u-STA-shun) tube, the narrow channel connecting the ear with the back of the throat. As the link between the middle ear and the throat, the eustachian tube has two roles. It helps drain normal, cleansing moisture from the middle ear. It also controls air pressure inside the middle ear chamber. When you swallow, the eustachian tube opens. This balances the air pressure in the middle ear with the pressure outside the eardrum. In infants and young children, the eustachian tube is short and almost level with the ear canal. By about age 7, however, the eustachian tube has become longer and steeper. This improves how well it works. The eardrum and middle ear are important to normal hearing. Together, they pass sound from the outer to the inner ear. When sound from the outer ear hits a flexible eardrum, the eardrum vibrates. The small bones in the middle ear pick up these vibrations and pass them along to the inner ear. There, the vibrations become electrical signals, which are sent along nerve pathways to the brain.

http://www.einstein.edu/einsteinhealthtopic/?languagecode=es&healthTopicId=-1&healthTopicName=HealthSheets&articleId=84614&articleTypeId=3

Out of all of the amazing wonders and phenomenon in space, one of the most amazing and perplexing is the black hole. A black hole has incredibly strong gravitational pull capturing everything from rock and gas to light. Nothing escapes a black hole. Scientists have studied black holes for years in an attempt to understand how they are created and how they work. Astronomers have announced that they have made the first reliable measurement of a supermassive black hole’s spin. The black hole the scientists studied is at the center of a spiral galaxy called NGC 1365. According to the scientists, this gigantic black hole is spinning at about 84% of the speed that Einstein’s general theory of relativity will allow. The discovery shows that some giant supermassive black hole’s are spinning very quickly. This discovery confirms studies in the past that had hypothesized black holes spun very quickly. Those previous studies were unable to confirm the speedy rotation of the target black hole. NGC 1365 is 56 million light-years away from Earth. The galaxy is located in the constellation Fornax, and the black hole in its center is more massive than several million suns. This particular black hole is also known to spew huge amounts of energy as it gobbles up gas and other objects nearby. The scientists used the high-energy light emitted by iron atoms consumed by the black hole to trace the motion of the flat, rotating accretion disk that circles the black hole.

http://www.slashgear.com/scientists-reliably-measure-a-supermassive-black-holes-spin-for-the-first-time-28271884/

Researchers have demonstrated tiny solar cells just billionths of a metre across that can repair themselves, extending their useful lifetime. The cells make use of proteins from the machinery of plants, turning sunlight into electric charges that can do work. The cells simply assemble themselves from a mixture of the proteins, minute tubes of carbon and other materials. The self-repairing mechanism, reported in Nature Chemistry, could lead to much longer-lasting solar cells. The design and improvement of solar cells is one of the most vibrant areas of science, in part because sunlight is far and away the planet's most abundant renewable energy source. More than that, nature has already proven that sunlight can be captured and turned into other forms of energy not only with extraordinary efficiency but also with a self-repair mechanism that counteracts the ravages of sunlight. "Sunlight, when it hits oxygen, is very damaging," explained Michael Strano, the Massachusetts Institute of Technology chemical engineer who led the research. "It's the reason why we age, and the reason why when we leave paper or plastic out in the sun, it fades." The destructive mixture of sunlight and oxygen, Professor Strano told BBC News, means that many of the best solar cells in the laboratory might not survive well when put into use. "There's a kind of a horse race among scientists around the world to make the highest efficiency cell, but very few people are asking what happens with that cell when you plug it in for a few hours or for a week or for months," he said.

http://tvnewslies.org/tvnl/index.php/news/science/15814-tiny-solar-cells-fix-themselves.html

September 19, 2011 The Central Valley of California is a fertile bed of over 350 diverse species of agricultural crops, some of the major cash crops being rice, grapes, cotton, and almonds. California’s agricultural industry makes up 15% of the entire nation’s crops and made a profit of $37 billion in the year 2009. Despite these numbers, desertification is an increasingly major problem. Between 1998 and 2000, 10,000 acres of farmland were lost every year in Central Valley from urbanization alone—this doesn’t account for the acres of fertile farmland lost due to overgrazing, climate change, or poor farming practices. Currently, California is losing 178 km2 of arable, fertile land each year. Southern California especially, being a very arid and drought-inclined region to begin with, has a problem with increasing salinity and compound minerals in the soil, caused by overdrawing ground water (United States Geological Survey). Desertification is not only the result of human activity. The UN Convention to Combat Desertification identifies the other major cause of desertification to be climatic variations—for example, erosion, drought and irregular rainfall, and violent winds. Essentially it renders the soil infertile, not only for planting and agriculture but for any organic life. Desertification occurs on a global scale, particularly through deforestation and drought. Areas around the Amazons, for example, have undergone desertification because the trees are being harvested for wood and cleared for farmland, and much of the space lies fallow. Similarly, in California, trees are cleared using the “slash and burn” method to open fields for cheap soybean and livestock cultivation. Desertification is a challenge for California because it is a desert environment supporting an increasingly large population on limited water imports. The situation becomes more dire when the effects of global warming are considered, which dramatically expedite desertification. Owens Valley, California, for example, became a desert when all of the natural water resources were diverted to Southern California for drinking water and crop irrigation. The San Joaquin Valley is a region that has undergone natural desertification due to climate change, a result of surface crusting, salinization and waterlogging problems. Most popularly considered solutions to desertification involve addressing problems of drought. Every 5 years in California a new Drought Contingency Plan (DCP) is released. In 2008, the last DCP, tactics included aggressive conservation, new groundwater and surface water storage facilities, and environmental restoration. GMOs also offer the possibility of growing crops that are resistant to drought, thus using less of the precious water resources to yield the same or greater amount of agriculture. Similarly, Air-to-Water harvesters are a new technology that essentially takes the humidity in the air and convert it to usable water. This can slow desertification significantly. From a more bottom-up perspective, education and conservation initiatives will also drastically reduce the human contribution to desertification. Programs teaching grey-water usage, water conservation, and the transformation of lawns into food forests can save a lot of water if it is implemented locally and broadly. Natural forests and wetlands need to be protected rather than cleared, farmed, and abandoned. While desertification is in and of itself a natural process, the human factors can and must be reduced, especially in California, if we are to live harmoniously with the land and reap the benefits of its yield. About the authors: Xueyou Wang and Kayla Duarte are undergraduate students in the USC Dana and David Dornsife College of Letters, Arts and Sciences.

http://dornsife.usc.edu/enst-320a/water-and-soil/?tag=human-activity

What You Need to Know A law is a scientific statement that is generally accepted as true. Newton's third law of motion says that for every action force there is an equal and opposite reaction force. How Does Newton's Third Law of Motion Work? Sir Isaac Newton (1642–1727) realized that if one object applies a force on another, the second object applies the same amount of force on the first object but in an opposite direction. Since these forces are not opposed, they are said to be unbalanced. Unbalanced forces cause the objects they push against to move in the direction of the force. The diagram on the left shows Newton's third law in action. When the boy hits the golf ball with the golf club, the ball is pushed forward, and the golf club is pushed backward. Two forces, A and B as shown, are needed to make these movements. Notice that the arrows for the forces are equal in size, are in opposite directions, and are on different objects. Force A is from the club hitting the ball, and force B is from the ball hitting the club. You can be sure that two forces are action-reaction pairs of forces if the reverse description of one force describes the other. In the figure, the identified action-reaction pair is A/B. The description of force A is "the club pushes against the ball," and the description of force B is "the ball pushes against the club." What Does This Have to Do with Creating a Balloon Rocket? If an inflated balloon with an open end is released, the balloon will fly through the air due to the unbalanced forces making up action/reaction pairs of forces. In the diagram below, one set of action/reaction forces (A/B) is shown. Force A is on the balloon due to the gas inside the balloon pushing on the wall of the balloon. Force B is on the gas due to the balloon pushing on the gas inside the balloon. Because of these unbalanced forces, the gas is pushed out of the balloon and the balloon is pushed forward. Action/reaction force pairs also make it possible for birds to fly. The wings of a bird push air downward (action). In turn, the air pushes the bird upward with an equal force (reaction). Real-Life Science Challenge How does a rocket fly through empty space? Like the balloon, real rockets do not move because the exhaust gas pushes against the ground or air surrounding the craft. Instead, a rocket moves because of the unbalanced forces in action/reaction pairs. That's why rockets move in space where there is nothing for the exhaust gas to push against. Now, start experimenting with your balloon rocket. What kind of balloon rocket goes the fastest? What kind goes the farthest? What effect do balloon size and shape have on the rocket's flight? - Design a way to control the direction of the balloon rocket's flight. For example, you could run a string through a straw, then stretch the string between two objects, such as chairs. Tape the balloon to the straw. - What else can make your balloon go farther, faster, or straighter? Warning is hereby given that not all Project Ideas are appropriate for all individuals or in all circumstances. Implementation of any Science Project Idea should be undertaken only in appropriate settings and with appropriate parental or other supervision. Reading and following the safety precautions of all materials used in a project is the sole responsibility of each individual. For further information, consult your state’s handbook of Science Safety.

http://www.education.com/science-fair/article/create-balloon-rocket/

The Digestive System - Design: parts of the digestive system The digestive system may be broken into two parts: a long, winding, muscular tube accompanied by accessory digestive organs and glands. That open-ended tube, known as the alimentary canal or digestive tract, is composed of various organs. These organs are, in order, the mouth, pharynx, esophagus, stomach, small intestine, and large intestine. The rectum and anus form the end of the large intestine. The accessory digestive organs and glands that help in the digestive process include the tongue, teeth, salivary glands, pancreas, liver, and gall bladder. The walls of the alimentary canal from the esophagus through the large intestine are made up of four tissue layers. The innermost layer is the mucosa, coated with mucus. This protects the alimentary canal from chemicals and enzymes (proteins that speed up the rate of chemical reactions) that break down food and from germs and parasites that might be in that food. Around the mucosa is the submucosa, which contains blood vessels, nerves, and lymph vessels. Wrapped around the submucosa are two layers of muscles that help move food along the canal. The outermost layer, the serosa, is moist, fibrous tissue that protects the alimentary canal and helps it move against the surrounding organs in the body. Food enters the body through the mouth, or oral cavity. The lips form and protect the opening of the mouth, the cheeks form its sides, the tongue forms its floor, and the hard and soft palates form its roof. The hard palate is at the front; the soft palate is in the rear. Attached to the soft palate is a fleshy, fingerlike projection called the uvula (from the Latin word meaning "little grape"). Two U-shaped rows of teeth line the mouth—one above and one below. Three pair of salivary glands open at various points into the mouth. - Alimentary canal (al-i-MEN-tah-ree ka-NAL): - Also known as the digestive tract, the series of muscular structures through which food passes while being converted to nutrients and waste products; includes the oral cavity, pharynx, esophagus, stomach, large intestine, and small intestine. - Amylase (am-i-LACE): - Any of various digestive enzymes that convert starches to sugars. - Appendix (ah-PEN-dix): - Small, apparently useless organ extending from the cecum. - Greenish yellow liquid produced by the liver that neutralizes acids and emulsifies fats in the duodenum. - Bolus (BO-lus): - Rounded mass of food prepared by the mouth for swallowing. - Cecum (SEE-kum): - Blind pouch at the beginning of the large intestine. - Chyle (KILE): - Thick, whitish liquid consisting of lymph and tiny fat globules absorbed from the small intestine during digestion. - Chyme (KIME): - Soupylike mixture of partially digested food and stomach secretions. - Colon (KOH-lun): - Largest region of the large intestine, divided into four sections: ascending, transverse, descending, and sigmoid (colon is sometimes used to describe the entire large intestine). - Colostomy (kuh-LAS-tuh-mee): - Surgical procedure where a portion of the large intestine is brought through the abdominal wall and attached to a bag to collect feces. - Defecation (def-e-KAY-shun): - Elimination of feces from the large intestine through the anus. - Dentin (DEN-tin): - Bonelike material underneath the enamel of teeth, forming the main part. - Duodenum (doo-o-DEE-num or doo-AH-de-num): - First section of the small intestine. - Emulsify (e-MULL-si-fie): - To break down large fat globules into smaller droplets that stay suspended in water. - Enamel (e-NAM-el): - Whitish, hard, glossy outer layer of teeth. - Enzymes (EN-zimes): - Proteins that speed up the rate of chemical reactions. - Epiglottis (ep-i-GLAH-tis): - Flaplike piece of tissue at the top of the larynx that covers its opening when swallowing is occurring. - Esophagus (e-SOF-ah-gus): - Muscular tube connecting the pharynx and stomach. - Feces (FEE-seez): - Solid body wastes formed in the large intestine. - Flatus (FLAY-tus): - Gas generated by bacteria in the large intestine. - Gastric juice (GAS-trick JOOSE): - Secretion of the gastric glands of the stomach, containing hydrochloric acid, pepsin, and mucus. - Ileocecal valve (ill-ee-oh-SEE-kal VALV): - Sphincter or ring of muscule that controls the flow of chyme from the ileum to the large intestine. - Ileum (ILL-ee-um): - Final section of the small intestine. - Jejunum (je-JOO-num): - Middle section of the small intestine. - Lacteals (LAK-tee-als): - Specialized lymph capillaries in the villi of the small intestine. - Larynx (LAR-ingks): - Organ between the pharynx and trachea that contains the vocal cords. - Lipase (LIE-pace): - Digestive enzyme that converts lipids (fats) into fatty acids. - Lower esophageal sphincter (LOW-er i-sof-ah-GEE-alSFINGK-ter): - Strong ring of muscle at the base of the esophagus that contracts to prevent stomach contents from moving back into the esophagus. - Palate (PAL-uht): - Roof of the mouth, divided into hard and soft portions, that separates the mouth from the nasal cavities. - Papillae (pah-PILL-ee): - Small projections on the upper surface of the tongue that contain taste buds. - Peristalsis (per-i-STALL-sis): - Series of wavelike muscular contractions that move material in one direction through a hollow organ. - Pharynx (FAR-inks): - Short, muscular tube extending from the mouth and nasal cavities to the trachea and esophagus. - Plaque (PLACK): - Sticky, whitish film on teeth formed by a protein in saliva and sugary substances in the mouth. - Pyloric sphincter (pie-LOR-ick SFINGK-ter): - Strong ring of muscle at the junction of the stomach and the small intestine that regulates the flow of material between them. - Rugae (ROO-jee): - Folds of the inner mucous membrane of organs, such as the stomach, that allow those organs to expand. - Trypsin (TRIP-sin): - Digestive enzyme that converts proteins into amino acids; inactive form is trypsinogen. - Uvula (U-vue-lah): - Fleshy projection hanging from the soft palate that raises to close off the nasal passages during swallowing. - Vestigial organ (ves-TIJ-ee-al OR-gan): - Organ that is reduced in size and function when compared with that of evolutionary ancestors. - Villi (VILL-eye): - Tiny, fingerlike projections on the inner lining of the small intestine that increase the rate of nutrient absorption by greatly increasing the intestine's surface area. THE TONGUE. The muscular tongue is attached to the base of the mouth by a fold of mucous membrane. On the upper surface of the tongue are small projections called papillae, many of which contain taste buds (for a discussion of taste, see chapter 12). Most of the tongue lies within the mouth, but its base extends into the pharynx. Located at the base of the tongue are the lingual tonsils, small masses of lymphatic tissue that serve to prevent infection. TEETH. Humans have two sets of teeth: deciduous and permanent. The deciduous teeth (also known as baby or milk teeth) start to erupt through the gums in the mouth when a child is about six months old. By the age of two, the full set of twenty teeth has developed. Between the ages of six and twelve, the roots of these teeth are reabsorbed into the body and the teeth begin to fall out. They are quickly replaced by the thirty-two permanent adult teeth. (The third molars, the wisdom teeth, may not erupt because of inadequate space in the jaw. In such cases, they become impacted or embedded in the jawbone and must be removed surgically.) Teeth are classified according to shape and function. Incisors, the chisel-shaped front teeth, are used for cutting. Cuspids or canines, the pointed teeth next to the incisors, are used for tearing or piercing. Bicuspids (or premolars) and molars, the back teeth with flattened tops and rounded, raised tips, are used for grinding. Each tooth consists of two major portions: the crown and the root. The crown is the exposed part of the tooth above the gum line; the root is enclosed in a socket in the jaw. The outermost layer of the crown is the whitish enamel. Made mainly of calcium, enamel is the hardest substance in the body. Underneath the enamel is a yellowish, bonelike material called dentin. It forms the bulk of the tooth. Within the dentin is the pulp cavity, which receives blood vessels and nerves through a narrow root canal at the base of the tooth. THE SALIVARY GLANDS. Three pair of salivary glands produce saliva on a continuous basis to keep the mouth and throat moist. The largest pair, the parotid glands, are located just below and in front of the ears. The next largest pair, the submaxillary or submandibular glands, are located in the lower jaw. The smallest pair, the sublingual glands, are located under the tongue. Ivan Petrovich Pavlov (1849–1936) was a Russian physiologist (a person who studies the physical and chemical processes of living organisms) who conducted pioneering research into the digestive activities of mammals. His now-famous experiments with a dog ("Pavlov's dog") to show how the central nervous system affects digestion earned him the Nobel Prize for Medicine or Physiology in 1904. Interested in the actions of digestion and gland secretion, Pavlov set up an ingenious experiment. In a laboratory, he severed a dog's throat (Pavlov was a skillful surgeon and the animal was unharmed). When the dog ate food, the food dropped out of the animal's throat before reaching its stomach. Through this simulated feeding, Pavlov discovered that the sight, smell, and swallowing of food was enough to cause the secretion of gastric juice. He demonstrated that the stimulation of the vagus nerve (one of the major nerves of the brain) influences the actions of the gastric glands. In another famous study, Pavlov set out to determine whether he could turn unconditioned (naturally occurring) reflexes or responses of the central nervous system into conditioned (learned) reflexes. He had noticed that laboratory dogs would sometimes salivate merely at the approach of lab assistants who fed them. Pavlov then decided to ring a bell each time a dog was given food. After a while, he rang the bell without feeding the dog. He discovered that the dog salivated at the sound of the bell, even though food was not present. Through this experiment, Pavlov demonstrated that unconditioned reflexes (salivation and gastric activity) could become conditioned reflexes that were triggered by a stimulus (the bell) that previously had no connection with the event (eating). Ducts or tiny tubes carry saliva from these glands into the mouth. Ducts from the parotid glands open into the upper portion of the mouth; ducts from the submaxillary and sublingual glands open into the mouth beneath the tongue. The salivary glands are controlled by the autonomic nervous system, a division of the nervous system that functions involuntarily (meaning the processes it controls occur without conscious effort on the part of an individual). The glands produce between 1.1 and 1.6 quarts (1 and 1.5 liters) of saliva each day. Although the flow is continuous, the amount varies. Food (or anything else) in the mouth increases the amount produced. Even the sight or smell of food will increase the flow. Saliva is mostly water (about 99 percent), with waste products, antibodies, and enzymes making up the small remaining portion. At mealtimes, saliva contains large quantities of digestive enzymes that help break down food. Saliva also controls the temperature of food (cooling it down or warming it up), cleans surfaces in the mouth, and kills certain bacteria present in the mouth. The pharynx, or throat, is a short, muscular tube extending about 5 inches (12.7 centimeters) from the mouth and nasal cavities to the esophagus and trachea (windpipe). It serves two separate systems: the digestive system (by allowing the passage of solid food and liquids) and the respiratory system (by allowing the passage of air). The esophagus, sometimes referred to as the gullet, is the muscular tube connecting the pharynx and stomach. It is approximately 10 inches (25 centimeters) in length and 1 inch (2.5 centimeters) in diameter. In the thorax (area of the body between the neck and the abdomen), the esophagus lies behind the trachea. At the base of the esophagus, where it connects with the stomach, is a strong ring of muscle called the lower esophageal sphincter. Normally, this circular muscle is contracted, preventing contents in the stomach from moving back into the esophagus. The stomach is located on the left side of the abdominal cavity just under the diaphragm (a membrane of muscle separating the chest cavity from the abdominal cavity). When empty, the stomach is shaped like the letter J and its inner walls are drawn up into long, soft folds called rugae. When the stomach expands, the rugae flatten out and disappear. This allows the average adult stomach to hold as much as 1.6 quarts (1.5 liters) of material. The dome-shaped portion of the stomach to the left of the lower esophageal sphincter is the fundus. The large central portion of the stomach is the body. The part of the stomach connected to the small intestine (the curve of the J) is the pylorus. The pyloric sphincter is a muscular ring that regulates the flow of material from the stomach into the small intestine by variously opening and contracting. That material, a soupylike mixture of partially digested food and stomach secretions, is called chyme. The stomach wall contains three layers of smooth muscle. These layers contract in a regular rhythm—usually three contractions per minute—to mix and churn stomach contents. Mucous membrane lines the stomach. Mucus, the thick, gooey liquid produced by the cells of that membrane, helps protect the stomach from its own secretions. Those secretions—acids and enzymes—enter the stomach through millions of shallow pits that open onto the surface of the inner stomach. Called gastric pits, these openings lead to gastric glands, which secrete about 1.6 quarts (1.5 liters) of gastric juice each day. Gastric juice contains hydrochloric acid and pepsin. Pepsin is an enzyme that breaks down proteins; hydrochloric acid kills microorganisms and breaks down cell walls and connective tissue in food. The acid is strong enough to burn a hole in carpet, yet the mucus produced by the mucous membrane prevents it from dissolving the lining of the stomach. Even so, the cells of the mucous membrane wear out quickly: the entire stomach lining is replaced every three days. Mucus also aids in digestion by keeping food moist. William Beaumont (1785–1853) was an American surgeon who served as an army surgeon during the War of 1812 (1812–15) and at various posts after the war. It was at one of these posts that he saw what perhaps no one before him had seen: the inner workings of the stomach. In 1882, while serving at Fort Mackinac in northern Michigan, Beaumont was presented with a patient named Alexis St. Martin. The French Canadian trapper, only nineteen at the time, has been accidently shot in the stomach. The bullet had torn a deep chunk out of the left side of St. Martin's lower chest. At first, no one thought he would survive, but amazingly he did. However, his wound never completely healed, leaving a 1 inch-wide (2.5 centimeter-wide) opening. This opening allowed Beaumont to put his finger all the way into St. Martin's stomach. Beaumont decided to take advantage of opening into St. Martin's side to study human digestion. He started by taking small chunks of food, tying them to a string, then inserting them directly into the young man's stomach. At irregular intervals, he pulled the food out to observe the varying actions of digestion. Later, using a hand-held lens, Beaumont peered into St. Martin's stomach. He observed how the human stomach behaved at various stages of digestion and under differing circumstances. Beaumont conducted almost 240 experiments on St. Martin. In 1833, he published his findings in Experiments and Observations on the Gastric Juice and the Physiology of Digestion , a book that provided invaluable information on the digestive process. The small intestine The small intestine is the body's major digestive organ. Looped and coiled within the abdominal cavity, it extends about 20 feet (6 meters) from the stomach to the large intestine. At its junction with the stomach, it measures about 1.5 inches (4 centimeters) in diameter. By the time it meets the large intestine, its diameter has been reduced to 1 inch (2.5 centimeters). Although much longer than the large intestine, the small intestine is called "small" because its overall diameter is smaller. The small intestine is divided into three regions or sections. The first section, the duodenum, is the initial 10 inches (25 centimeters) closest to the stomach. Chyme from the stomach and secretions from the pancreas and liver empty into this section. The middle section, the jejunum, measures about 8.2 feet (2.5 meters) in length. Digestion and the absorption of nutrients occurs mainly in the jejunum. The final section, the ileum, is also the longest, measuring about 11 feet (3.4 meters) in length. The ileum ends at the ileocecal valve, a sphincter that controls the flow of chyme from the ileum to the large intestine. The inner lining of the small intestine is covered with tiny, fingerlike projections called villi (giving it an appearance much like the nap of a plush, soft towel). The villi greatly increase the intestinal surface area available for absorbing digested material. Within each villus (singular for villi) are blood capillaries and a lymph capillary called a lacteal. Digested food molecules are absorbed through the walls of the villus into both the capillaries and the lacteal. At the bases of the villi are openings of intestinal glands, which secrete a watery intestinal juice. This juice contains digestive enzymes that convert food materials into simple nutrients the body can readily use. On average, about 2 quarts (1.8 liters) of intestinal juice are secreted into the small intestine each day. As with the lining of the stomach, a coating of mucus helps protect the lining of the small intestine. Yet again, the digestive enzymes prove too strong for the delicate cells of that lining. They wear out and are replaced about every two days. The large intestine Extending from the end of the small intestine to the anus, the large intestine measures about 5 feet (1.5 meters) in length and 3 inches (7.5 centimeters) in diameter. It almost completely frames the small intestine. The large intestine is divided into three major regions: the cecum, colon, and rectum. Cecum comes from the Latin word caecum , meaning "blind." Shaped like a rounded pouch, the cecum lies immediately below the area where the ileum empties into the large intestine. Attached to the cecum is the slender, fingerlike appendix, which measures about 3.5 inches (9 centimeters) in the average adult. Composed of lymphatic tissue, the appendix seems to have no function in present-day humans. For that reason, scientists refer to it as a vestigial organ (an organ that is reduced in size and function when compared with that of evolutionary ancestors). Sometimes used to describe the entire large intestine, the colon is actually the organ's main part. It is divided into four sections: ascending, transverse, descending, and sigmoid. The ascending colon travels from the cecum up the right side of the abdominal cavity until it reaches the liver. It then makes a turn, becoming the transverse colon, which travels horizontally across the abdominal cavity. Near the spleen on the left side, it turns down to form the descending colon. At about where it enters the pelvis, it becomes the S-shaped sigmoid colon. After curving and recurving, the sigmoid colon empties into the rectum, a fairly straight, 6-inch (15-centimeter) tube ending at the anus, the opening to the outside. Two sphincters (rings of muscle) control the opening and closing of the anus. Roughly 1.6 quarts (1.5 liters) of watery material enters the large intestine each day. No digestion takes place in the large intestine, only the reabsorption or recovery of water. Mucus produced by the cells in the lining of the large intestine help move the waste material along. As more and more water is removed from that material, it becomes compacted into soft masses called feces. Feces are composed of water, cellulose and other indigestible material, and dead and living bacteria. The remnants of worn red blood cells gives feces their brown color. Only about 3 to 7 ounces (85 to 200 grams) of solid fecal material remains after the large intestine has recovered most of the water. That material is then eliminated through the anus, a process called defecation. The pancreas is a soft, pink, triangular-shaped gland that measures about 6 inches (15 centimeters) in length. It lies behind the stomach, extending from the curve of the duodenum to the spleen. While a part of the digestive system, the pancreas is also a part of the endocrine system, producing the hormones insulin and glucagon (for a further discussion of this process, see chapter 3). Primarily a digestive organ, the pancreas produces pancreatic juice that helps break down all three types of complex food molecules in the small intestine. The enzymes contained in that juice include pancreatic amylase, pancreatic lipase, and trypsinogen. Amylase breaks down starches into simple sugars, such as maltose (malt sugar). Lipase breaks down fats into simpler fatty acids and glycerol (an alcohol). Trypsinogen is the inactive form of the enzyme trypsin, which breaks down proteins into amino acids. Trypsin is so powerful that if produced in the pancreas, it would digest the organ itself. To prevent this, the pancreas produces trypsinogen, which is then changed in the duodenum to its active form. Pancreatic juice is collected from all parts of the pancreas through microscopic ducts. These ducts merge to form larger ducts, which eventually combine to form the main pancreatic duct. This duct, which runs the length of the pancreas, then transports pancreatic juice to the duodenum of the small intestine. The largest glandular organ in the body, the liver weighs between 3 and 4 pounds (1.4 and 1.8 kilograms). It lies on the right side of the abdominal cavity just beneath the diaphragm. In this position, it overlies and almost completely covers the stomach. Deep reddish brown in color, the liver is divided into four unequal lobes: two large right and left lobes and two smaller lobes visible only from the back. The liver is an extremely important organ. Scientists have discovered that it performs over 200 different functions in the body. Among its many functions are processing nutrients, making plasma proteins and blood-clotting chemicals, detoxifying (transforming into less harmful substances) alcohol and drugs, storing vitamins and iron, and producing cholesterol. One of the liver's main digestive functions is the production of bile. A watery, greenish yellow liquid, bile consists mostly of water, bile salts, cholesterol, and assorted lipids or fats. Liver cells produce roughly 1 quart (1 liter) of bile each day. Bile leaves the liver through the common hepatic duct. This duct unites with the cystic duct from the gall bladder to form the common bile duct, which delivers bile to the duodenum. In the small intestine, bile salts emulsify fats, breaking them down from large globules into smaller droplets that stay suspended in the watery fluid in the small intestine. Bile salts are not enzymes and, therefore, do not digest fats. By breaking down the fats in to smaller units, bile salts aid the fatdigesting enzymes present in the small intestine. The gall bladder The gall bladder is a small, pouchlike, green organ located on the undersurface of the right lobe of the liver. It measures 3 to 4 inches (7.6 to 10 centimeters) in length. The gall bladder's function is to store bile, of which it can hold about 1.2 to 1.7 ounces (35 to 50 milliliters). The liver continuously produces bile. When digestion is not occurring, bile backs up the cystic duct and enters the gall bladder. While holding the bile, the gall bladder removes water from it, making it more concentrated. When fatty food enters the duodenum once again, the gall bladder is stimulated to contract and spurt out the stored bile.

http://www.faqs.org/health/Body-by-Design-V1/The-Digestive-System-Design-parts-of-the-digestive-system.html

Kentucky's Underground Railroad: Passages to Freedom fills in a missing piece of Kentucky history and supports inquiry in a variety of subject areas. Teachers in grades 4-12 can use the entire 60-minute program, or segments, to illustrate Kentucky's role in the story of slavery, abolitionism, and the underground railroad. Several Kentucky educators have reviewed the documentary and written letters to colleagues suggesting uses with students at the elementary, middle and high school levels. Letters to Colleagues An overview of some of the themes in the documentary can be used to develop cross curricular units. Individual sections of the documentary can be adapted to fit specific grade levels. Plus, a behind-the-scenes look at the documentary provides ideas to incorporate technology with content and offers insights into careers in television production. Visit the KDE companion Web site (www.kde.state.ky.us/) for more teaching resources and related educational Web sites. The Kentucky Department of Education Core Content for Assessment site will be useful to teachers in all curriculum areas. KDE Core Content Kentucky's Learning Goals and Academic Expectations Kentucky's Underground Railroad-Passage to Freedom connects to the following goals in social studies and arts and humanities. Arts and Humanities Documentary Theme Overview Kentucky's Underground Railroad-Passage to Freedom illustrates the influence of Kentucky geography, history, economy, politics, and culture in the context of a much larger social, political and moral struggle in American history. By hearing interviews with real people telling their family stories, students will realize that history lessons are not always found in books. They will also realize the importance of corroboration when history was not recorded and activities were conducted secretly. Humanities themes in this documentary are: This story of rebellion and sacrifice in resistance to human bondage and cruelty has largely been undocumented. The costs of being discovered after a failed escape were death, imprisonment, whipping, or sale "down the river" so secrets of successful escapes were well hidden. Help for escaping slaves could come from many sources including other slaves, free blacks (some never held in bondage), Native Americans, whites acting alone, or whites acting as conductors along the Underground Railroad. Similarly, there were patrollers acting alone or with others who tried to foil escapes and return slaves to their masters. Secret codes of communication to aid escapees were embedded in music, quilts, and signals. The stars, particularly the North Star, were guides for runaways. Though hidden passageways, stairs, and rooms are frequently rumored throughout Kentucky, few are actually documented as sites to aid the escape of slaves. Importance of documentation and preservation of local history Many of the stories represented in the documentary have been passed down for generations in oral histories, some with physical evidence that the stories were true. Others are well documented and represent some of the bravest stories of Kentuckians on record. Kentucky still has the opportunity to document and preserve what remains untold of this story. Teachers and students must first incorporate what is documented into the curriculum and then encourage students at all levels to research their own family and local community histories and use various forms of communication to share with others.

http://www.ket.org/underground/resources/

The CERN Large Hadron Collider (LHC) dipole magnets are designed to operate at magnetic fields of about eight tesla, or 8 T These LHC magnets use cables made of superconducting niobium titanium (NbTi), and for five years during its construction the LHC contracted for more than 28 percent of the world’s niobium titanium wire production, with significant quantities of NbTi also used in the magnets for the LHC’s giant experiments. Berkeley labs is working with the high-temperature superconductor Bi-2212 (bismuth strontium calcium copper oxide) to develop stronger magnets. In experimental situations Bi-2212 has generated fields of 25 T (tesla) and could go much higher. But like many high-temperature superconductors Bi-2212 is not a metal alloy but a ceramic, virtually as brittle as a china plate. One of the things that makes Bi-2212 promising is that it is now available in the form of round wires. “The wires are essentially tubes filled with tiny particles of ground-up Bi-2212 in a silver matrix,” Godeke explains. “While the individual particles are superconducting, the wires aren’t – and can’t be, until they’ve been heat treated so the individual particles melt and grow new textured crystals upon cooling – thus welding all of the material together in the right orientation.” Orientation is important because Bi-2212 has a layered crystalline structure in which current flows only through two-dimensional planes of copper and oxygen atoms. Out of the plane, current can’t penetrate the intervening layers of other atoms, so the copper-oxygen planes must line up if current is to move without resistance from one Bi-2212 particle to the next. In a coil fabrication process called “wind and react,” the wires are first assembled into flat cables and the cables are wound into coils. The entire coil is then heated to 888 degrees Celsius (888 C) in a pure oxygen environment. During the “partial melt” stage of the reaction, the temperature of the coil has to be controlled to within a single degree. It’s held at 888 C for one hour and then slowly cooled. Silver is the only practical matrix material that allows the wires to “breathe” oxygen during the reaction and align their Bi-2212 grains. Unfortunately 888 C is near the melting point of silver, and during the process the silver may become too soft to resist high stress, which will come from the high magnetic fields themselves: the tremendous forces they generate will do their best to blow the coils apart. So far, attempts to process coils have often resulted in damage to the wires, with resultant Bi-2212 current leakage, local hot spots, and other problems. The goal of the program to develop Bi-2212 for high-field magnets is to improve the entire suite of wire, cable, coil making, and magnet construction technologies. We need to improve current density by a factor of three or four. The Large Dipole test magnet, called LD1, will be used to test magnets and wires. The test magnet will be based on niobium-tin technology and will exert a field of up to 15 T across the height of the aperture. Inside the aperture, two cable samples will be arranged back to back, running current in opposite directions to minimize the forces generated by interaction between the sample and the external field applied by LD1. The magnets are being developed to make the highest-energy colliders possible. The new technology will benefit many other fields as well, from undulators for next-generation light sources to more compact medical devices. If you liked this article, please give it a quick review on Reddit, or StumbleUpon. Thanks How to Make Money

http://nextbigfuture.com/2010/09/berkeley-labs-developing-stronger.html

After planning your action, comes the time to manage really your classroom. Here are some tips to guide you. ICT and differentiated instruction More and more students with disabilities are taught in an inclusive setting.. Fortunately, ICT is an excellent tool to practice differentiated instruction. A growing body of research confirms that when it comes to learning, one size doesn't fit all. Differentiated instruction is based on the idea of adapting the teaching strategy and content to the needs, speed and learning styles of each student. Adapting the context: Classroom layout must be accessible for all students. For example, teacher must make sure that there is enough space to accommodate students using a wheelchair. Adapting the procedures: It consists in introducing modifications on teaching strategies or the tasks given to the students. As an example, by creating a code of colors for the tasks (e.g., yellow circles for reading tasks and a blue squares for exercises), you will create inclusive teaching materials that may assist an ESL (English as a Second Language) student as well as others who learn best by using visual aids. Adapting the content: It is important to select flexible technologies that offer options in adapting the content to the students’ needs (for e.g. audio version of texts (Link to appropriate section in Alternative Formats)). Keep in mind to ensure that technology use (or the application of its flexibility options) does not interfere with other students in the class (disturbing noise, complicate interface because of the multiplication of the flexibility options, etc.) Adapting the evaluation: Technology can help you to adapt the output of learning of students, inclusive of those with disabilities. For example, the use of word processing is an excellent strategy to accommodate students who experience difficulty writing. 1. Technology isn't the spotlight – technology should support and encourage the use, the development and the quality educational content 2. Technology use must be adapted to the needs and characteristics of the students 3. Technology use must take into account the content of the individualized educational plan (IEP) of the student and the learning objectives it contains Take Action relates to the following sections of the Accessibility for Ontarians with Disabilities Act (AODA) Integrated Accessibility Regulation: Accessible Formats and Communication Supports Training on Accessible Course Delivery and Instruction 1. The Take Action section describes how educators may develop teaching strategies and content that support the various learning styles of students. Learning about who uses alternative formats and how these are used assists educators with integrating inclusive approaches to teach, communicate and share information. 2. People interact, learn and communicate in diverse ways. Learning opportunities are increased when flexible ways of engaging with learning materials are provided. Considering how people communicate is important for knowledge to be exchanged. Alternative formats take into account diverse ways of exchanging information. 3. The AODA legislates that educators, teachers and staff are to learn about accessible course delivery and instruction and be knowledgeable at interacting and communicating with people with disabilities who may use alternative formats.

http://snow.idrc.ocad.ca/node/119

Liberation of Slaves and Colonization in Liberia in 1825 In 1825, Christian groups and colonization societies in America advocated for freed African Americans to colonize land in Liberia as an alternative to emancipate slaves in America. Religious groups expressed their sentiments in publication including the Christian Register, which published an article in an issued dated September 3, 1825. The article argued for the transportation of slaves to their homeland in Liberia because it would be beneficial to the discriminated African American race for a plethora of reasons. The publication argued slaves would regain freedoms denied to them in America as well as allow them to establish their own government the way they desired. Freed slaves would be allowed to cultivate the land of Liberia by utilizing certain technology the United States would introduce to the country. Colonization societies believed light American presence would allow former slaves to excel in the agricultural, economic and political realms of their society because freed slaves in Liberia would have ties with one of the most technologically advanced nations of the era. The Christian Register boldly emphasized the emancipation of enslaved peoples as morally correct, and a resolution was desperately needed to appease slaveholders and abolitionists alike. Sending slaves to Liberia seemed to offer the best mode of compromise. The issue of whether the emancipation of enslaved African Americans was necessary during the nineteenth century played a crucial role in the development of beliefs in certain groups, such as the American Colonization Society and the Pennsylvania Colonization Society. The American Colonization Society, the Pennsylvania Colonization Society, and the Christian Register advocated that the sending of freed slaves would be beneficial to enslaved African Americans. However, after reanalyzing the efforts of pro-colonization societies and publication, historians of the 21st century have come to understand that colonization was in response to the threat of freed African Americans if emancipation legislation was passed in the United States. Abolitionists offered information, which discredited the efforts of colonization societies by arguing reports from Liberia were deplorable. Slaves were rarely sent to their homeland and conditions in Liberia were treacherous. There was no possible way freed slaves could colonize Liberia the way pro-colonization supports had argued they could in the nineteenth century. Abolitionists were capable of discrediting the ideologies of colonization societies because it was based on the belief of negrophobia. In fear of how African Americans would respond if they were finally emancipated from slavery, colonization societies believed by sending them to Liberia, they would not have to deal with the potential hazardous outcomes of emancipation. - "Moral and Religious: Advantages of Colonization in Africa," Christian Register, September 3, 1825. - Burin, Eric, "Rethinking Northern White Support for the African Colonization Movement: The Pennsylvania Colonization Society as an Agent of Emancipation," Pennsylvania Magazine of History and Biography Vol CXXVII No 2 (April 2003): 197-229.

http://historyengine.richmond.edu/episodes/view/4449

The following HTML text is provided to enhance online readability. Many aspects of typography translate only awkwardly to HTML. Please use the page image as the authoritative form to ensure accuracy. Adding + It Up: Helping Children Learn Mathematics teachers should know. Many of these ideas are treated in more detail in textbooks intended for prospective elementary school teachers. A major theme of the chapter is that numbers are ideas—abstractions that apply to a broad range of real and imagined situations. Operations on numbers, such as addition and multiplication, are also abstractions. Yet in order to communicate about numbers and operations, people need representations—something physical, spoken, or written. And in order to carry out any of these operations, they need algorithms: step-by-step procedures for computation. The chapter closes with a discussion of the relationship between number and other important mathematical domains such as algebra, geometry, and probability. At first, school arithmetic is mostly concerned with the whole numbers: 0, 1, 2, 3, and so on.1 The child’s focus is on counting and on calculating— adding and subtracting, multiplying and dividing. Later, other numbers are introduced: negative numbers and rational numbers (fractions and mixed numbers, including finite decimals). Children expend considerable effort learning to calculate with these less intuitive kinds of numbers. Another theme in school mathematics is measurement, which forms a bridge between number and geometry. Mathematicians like to take a bird’s-eye view of the process of developing an understanding of number. Rather than take numbers a pair at a time and worry in detail about the mechanics of adding them or multiplying them, they like to think about whole classes of numbers at once and about the properties of addition (or of multiplication) as a way of combining pairs of numbers in the class. This view leads to the idea of a number system. A number system is a collection of numbers, together with some operations (which, for purposes of this discussion, will always be addition and multiplication), that combine pairs of numbers in the collection to make other numbers in the same collection. The main number systems of arithmetic are (a) the whole numbers, (b) the integers (i.e., the positive whole numbers, their negative counterparts, and zero), and (c) the rational numbers—positive and negative ratios of whole numbers, except for those ratios of a whole number and zero. Thinking in terms of number systems helps one clarify the basic ideas involved in arithmetic. This approach was an important mathematical discovery in the late nineteenth and early twentieth centuries. Some ideas of arithmetic are fairly subtle and cause problems for students, so it is useful to have a viewpoint from which the connections between ideas can be surveyed.

http://www.nap.edu/openbook.php?record_id=9822&page=72

Resistance and resistivity Components have resistance, materials have resistivity Resistance is a property of a particular component, like this resistor or this bulb. The symbol for resistance is R and it's measured in ohms (Ω). Practise past papers and get Better Grades! Visit our sister site www.good-at-maths.com for - PAST PAPERS - MODEL ANSWERS - HD VIDEO EXPLANATIONS Resistivity is a property of a material, like copper or plastic. The symbol for resistivity is ρ (the Greek small letter rho - r for rho, r for resistivity). It's measured in 'ohms metre'. You can see why below. When you talk about resistivity it doesn't matter how big or what shape the sample is. In this sense it's a bit like density because you can compare lead with wood without having to say which bit of lead or wood you're talking about. When you say copper is a better conductor than plastic you're really saying that lead has a lower resistivity than plastic. In other words you're making a general statement about lead and plastic, rather than particular components. Resistance, resistivity and size The resistance of piece of wire is big if it's long and thin. In other words the resistance of a particular piece of wire depends on - what metal it's made from (its resistivity, ρ) - how long it is (its length, l) - its cross-sectional area (its area, A, which decreases resistance the bigger it is) We can write R = ρl/A. If we rearrange this we get ρ = RA/l. This has units of Ωm2/m, which is just Ωm (said 'ohms metre'). Changes of resistance and resistivity Resistivity is normally taught in more advanced courses, which is why we tend to talk only about resistance. This becomes slightly clumsy because when you want to say 'the resistivity of a metal increases with temperature', you have to be more specific and say 'the resistance of a piece of metal increases with temperature'.

http://www.furryelephant.com/content/electricity/resistance-ohms-law/resistance-resistivity/

Greening of the Red Planet |Tweet|Greening of the Red Planet A hardy microbe from Earth might one day transform the barren ground of Mars into arable soil. January 26, 2001 -- Although Mars may once have been warm and wet, the Red Planet today is a frozen wasteland. Most scientists agree, it's highly unlikely that any living creature --even a microbe-- could survive for long on the surface of Mars. When the first humans travel there to explore the Red Planet up close, they will have to grow their food in airtight, heated greenhouses. The Martian atmosphere is far too cold and dry for edible plants to grow in the open air. But if humans ever hope to establish long-term colonies on their planetary neighbor, they will no doubt want to find a way to farm outdoors. Imre Friedmann has an idea of how they might take the first step. Above: Artists' James Graham and Kandis Elliot impression of a more habitable Mars. [more from ThinkQuest.org] Sign up for EXPRESS SCIENCE NEWS delivery Mars is covered by a layer of ground-up rock and fine dust, known as regolith. To convert regolith into soil, it will be necessary to add organic matter, much as organic farmers on Earth fertilize their soil by adding compost to it. On Earth, compost is made up primarily of decayed vegetable matter. Microorganisms play an important role in breaking down dead plants, recycling their nutrients back into the soil so that living plants can reuse them. But on Mars, says Friedmann, where there is no vegetation to decay, the dead bodies of the microorganisms themselves will provide the organic matter needed to build up the soil. The trick is finding the right microbe. "Among the organisms that are known today," says Friedmann, "Chroococcidiopsis is most suitable" for the task. Chroococcidiopsis is one of the most primitive cyanobacteria known. What makes it such a good candidate is its ability to survive in a wide range of extreme environments that are hostile to most other forms of life. Chroococcidiopsis has been found growing in hot springs, in hypersaline (high-salt) habitats, in a number of hot, arid deserts throughout the world, and in the frigid Ross Desert in Antarctica. Above: A photomicrograph of Chroococcidiopsis, enlarged 100 times. "Chroococcidiopsis is the constantly appearing organism in nearly all extreme environments," Friedmann points out, "at least extreme dry, extreme cold, and extremely salty environments. This is the one which always comes up." Moreover, where Chroococcidiopsis survives, it is often the only living thing that does. But it gladly gives up its dominance when conditions enable other, more complex forms of life to thrive. For clues on how to farm Chroococcidiopsis on Mars, Friedmann looks to its growth habits in arid regions on Earth. In desert environments, Chroococcidiopsis grows either inside porous rocks, or just underground, on the lower surfaces of translucent pebbles. Above: In many desert environments, Chroococcidiopsis grows on the undersides of transparent rocks, just below the surface. The pebbles provide an ideal microenvironment for Chroococcidiopsis in two ways. First, they trap moisture underneath them. Experiments have shown that small amounts of moisture can cling to the undersurfaces of rocks for weeks after their above-ground surfaces have dried out. Second, because the pebbles are translucent, they allow just enough light to reach the organisms to sustain growth. Friedmann envisions large farms where the bacteria are cultured on the underside of strips of glass that are treated to achieve the proper light-transmission characteristics. Mars today, however, is too cold for this technique to work effectively. Before even as hardy a microbe as Chroococcidiopsis could be farmed on Mars, the planet would have to be warmed up considerably, to just below the freezing point. Friedmann, pictured left, admits that his ideas about growing Chroococcidiopsis are, at this point, merely a thought experiment. "I don't think any of us alive today will see this happen," he muses. When the time does come to make Mars a more habitable place, "the technology will be so different that everything we plan today... will be ridiculously outdated." Friedmann fully expects that genetic engineering will eventually develop designer organisms to do the job. Even if Chroococcidiopsis is ultimately used as the basis, it will be a vastly improved version of today's microbe. The Physics and Biology of Making Mars Habitable -- Web page for the conference where Friedmann presented his research Bibliography on terraforming -- extensive list of publications about terraforming, compiled by Chris McKay of NASA's Astrobiology Institute The Terraforming Information Pages -- links to a variety of resources about terraforming Meet Conan the Bacterium -- Science@NASA article: a humble microbe could become "The Accidental (Space) Tourist" NASA Astrobiology Institute -- home page Join our growing list of subscribers - sign up for our express news delivery and you will receive a mail message every time we post a new story!!! |The Science and Technology Directorate at NASA's Marshall Space Flight Center sponsors the Science@NASA web sites. The mission of Science@NASA is to help the public understand how exciting NASA research is and to help NASA scientists fulfill their outreach responsibilities.| |For lesson plans and educational activities related to breaking science news, please visit Thursday's Classroom|| Production Editor: Dr. Tony Phillips Curator: Bryan Walls Media Relations: Steve Roy Responsible NASA official: John M. Horack

http://science1.nasa.gov/science-news/science-at-nasa/2001/ast26jan_1/

Hepatitis C is an infectious disease caused by the Hepatitis C virus, or HCV. It is spread through ... Read the full transcript » Hepatitis C is an infectious disease caused by the Hepatitis C virus, or HCV. It is spread through body fluids, most commonly blood. Once the virus enters the bloodstream, it travels to the liver and attaches to receptors on the surface of liver cells called hepatocytes. The virus's fatty or lipid coat merges with the cell's outer membrane, so that the virus's protein core is inside the cell. The coat on the core dissolves. This releases a strand of RNA, which carries viral genetic information. The virus takes over components of the cell called ribosomes, which produce an enzyme needed to reproduce the viral RNA. When there is enough of this enzyme, called RNA transcriptase, the viral RNA produces a copy of itself. This copy serves as a template, which is copied many times as the genetic material for new viruses. The viral RNA directs the ribosomes to produce proteins that form the virus's protein coat. They assemble around the RNA into new virus particles. The viruses travel to the membrane, which circle and release it with a new lipid coat. This process continues until the cell dies. Over time millions of liver cells can be destroyed. The immune system also attacks the infected cells. Both can damage the liver. About 75% of those infected develop chronic infection. At its mildest, there are no signs or symptoms. Serum enzymes - a blood test for liver disease - are normal, and liver biopsy shows only mild injury. Those with severe disease have symptoms, high blood levels of HCV RNA, elevated serum enzymes, and significant liver damage. At least 20% develop cirrhosis. Treatment options are available. See your doctor for further information.

http://www.healthline.com/vpvideo/what-is-hepatitis-c

On this day in 1915, the French National Assembly passes a law formally ceding the land that holds the British war cemeteries to Great Britain. The move ensured even as the war was being fought that its saddest and most sacred monuments would be forever protected. The law stated that the land was "the free gift of the French people for a perpetual resting place of those who are laid there." By the end of the war, it would apply to more than 1,200 cemeteries along the Western Front, the majority located near the battlefields in the Somme, Nord, and Pas-de-Calais regions. The Commonwealth War Graves Commission, established in 1917 by a British royal charter, supervised the construction of the cemeteries and their monuments, which were designed by some of the most prominent British architects of the day. The last monument was put in place in 1938. The French office of the commission is charged with the maintenance of these cemeteries; between 400 and 500 members of its staff tend the graves and the surrounding horticulture. In addition to the cemeteries, the Commonwealth War Graves Commission also tends to the numerous monuments that exist on the Western Front to commemorate the missing. One of the largest of these stands at Thiepval, on the Somme battlefield, and bears the names of 73,357 British and South African soldiers and officers who died there between July 1915 and March 1918 and whose final resting place is not known.

http://www.history.com/this-day-in-history/french-government-gives-land-for-british-war-cemeteries

Read the following sentences: - Mother cooks dinner. - Children play in the park. - Barking dogs seldom bite. In the sentences given above, the words in bold text are used to say something about a person or a thing. They say what a person or a thing does. These words are called verbs. Now read the following sentences. - We have two hands and two legs. - She is a good girl. Here the verbs have and is show what a person has or is. These words are also called verbs. Thus we have seen that a verb is a word which shows what a person or a thing is, has or does. The verb may also express what happens or is done to the person or thing. The thief was beaten. (Here the verb was beaten shows what happens to the thief.) A verb may consist of more than one word. Some verbs may consist of as many as four words. - It is raining. - It has been raining. - It rains. Transitive and intransitive verbs Verbs that take an object are called transitive verbs. - She heard a noise. (subject – she, verb – heard, object – a noise) - He saw a pigeon. (Subject – he, verb – saw, object – a pigeon) - The girl plucked the flower. (Subject – the girl, verb – plucked, object – the flower) - The master beat the dog. (subject – the master, verb – beat, object – the dog) Some verbs do not take an object after them. These are called intransitive verbs. Examples are: smile, sit, sleep, cry, laugh, dance etc. - The baby smiled. (Here the verb smiled is intransitive because it has no object.) - The child cried. (Here the verb cried is intransitive because it has no object.) - He sat on the bed. (Here the verb sat is intransitive because it has no object.) Note that most verbs can be used both transitively and intransitively.

http://www.englishgrammar.org/verbs/

The completion of the railroads to the West following the Civil War opened up vast areas of the region to settlement and economic development. White settlers from the East poured across the Mississippi to mine, farm, and ranch. African-American settlers also came West from the Deep South, convinced by promoters of all-black Western towns that prosperity could be found there. Chinese railroad workers further added to the diversity of the region's population. Settlement from the East transformed the Great Plains. The huge herds of American bison that roamed the plains were virtually wiped out, and farmers plowed the natural grasses to plant wheat and other crops. The cattle industry rose in importance as the railroad provided a practical means for getting the cattle to market. The loss of the bison and growth of white settlement drastically affected the lives of the Native Americans living in the West. In the conflicts that resulted, the American Indians, despite occasional victories, seemed doomed to defeat by the greater numbers of settlers and the military force of the U.S. government. By the 1880s, most American Indians had been confined to reservations, often in areas of the West that appeared least desirable to white settlers. The cowboy became the symbol for the West of the late 19th century, often depicted in popular culture as a glamorous or heroic figure. The stereotype of the heroic white cowboy is far from true, however. The first cowboys were Spanish vaqueros, who had introduced cattle to Mexico centuries earlier. Black cowboys also rode the range. Furthermore, the life of the cowboy was far from glamorous, involving long, hard hours of labor, poor living conditions, and economic hardship. The myth of the cowboy is only one of many myths that have shaped our views of the West in the late 19th century. Recently, some historians have turned away from the traditional view of the West as a frontier, a "meeting point between civilization and savagery" in the words of historian Frederick Jackson Turner. They have begun writing about the West as a crossroads of cultures, where various groups struggled for property, profit, and cultural dominance. Think about these differing views of the history of the West as you examine the documents in this collection. To find additional documents in American Memory on topics related to the West, use such keywords as West, ranching, Native Americans, and pioneers, or search using the names of states or cities in the West.

http://loc.gov/teachers/classroommaterials/presentationsandactivities/presentations/timeline/riseind/west/

The graph of a quadratic function of the form f(x) = a x 2 + b x + c is a parabola Properties of Graphs of Quadratic Functions. a) If a > 0, the parabola opens upward; if a < 0, the parabola opens downward. b) As | a| increases, the parabola becomes narrower; as | a| decreases, the parabola becomes wider. c) The lowest point of a parabola (when a > 0) or the highest point (when a < 0) is called the vertex. d) The domain of a quadratic function is R, because the graph extends indefinitely to the right and to the left. If (h, k) is the vertex of the parabola, then the range of the function is [k,+ ∞ ) when a > 0 and (- ∞, k] when a < 0. e) The graph of a quadratic function is symmetric with respect to a vertical line containing the vertex. This line is called the axis of symmetry. If (h, k) is the vertex of a parabola, then the equation of the axis of symmetry is x = h. How to calculate the vertex of a parabola 1. To determine the vertex of the graph of a quadratic function, f(x) = ax2+ bx + c, we can either do it: a) Completing the square to rewrite the function in the form f(x) = a(x – h)2 + k. The vertex is (h, k). b) Using the formula to find the x-coordinate of the vertex and then, the y- coordinate of the vertex can be determined by evaluating . The vertex is (x,y). How to calculate the intercepts of a parabola. To find the y-intercept of the parabola, find f(0); to find the x-intercepts, solve the quadratic equation ax2+ bx + c = 0.

http://www.emathematics.net/parabola.php

Three species of mammoths (genus Mammuthus) lived on the mainland of the United States at the end of the last Ice Age. These were the Columbian mammoth (M. columbi), Jefferson's mammoth (M. jeffersonii), and the woolly mammoth (M. primigenius). Of these, Jefferson's mammoth and the woolly mammoth have been identified from the midwestern U.S. Mammoths, mastodons and modern elephants, are members of the order Proboscidea. The mammoths are closely related to the living elephants, especially to the Asiatic elephant (Elephas maximus). As adults these late-occuring mammoths stood between about 3 and 3.7 meters (10-12 feet) at the shoulder and weighed between 5500 and 7300 kilograms (6-8 tons). The teeth of mammoths are quite distinctive. They are composed of a set of compressed enamel plates that are held together with cementum. These cemented plates make a very tall, strong, and wear-resistant tooth. After a tooth erupts from the gum cavity, the mammoth uses it in grinding coarse vegetation like grass. This use causes the tooth to develop a flat top with low enamel ridges where the plates have been worn. The tall structure of these hypsodont (shallow-rooted) teeth make them very resistant to wear. This is important because mammoths are thought to have been primarily grass-eaters. Grass is a very hard material to eat. It has small pieces of silica (a glass-like substance) in its leaves. These pieces of silica act like sandpaper grit and would wear away a less resistant tooth very quickly. Mammoths are frequently found as fossils in the midwestern U.S. Most often isolated teeth are found. Mammoth fossils are most common in areas that were covered by savannas, grasslands, or tundra during the last Ice Age. This map shows some of the important mammoth finds in the region. Approximately 1.5 to 1.8 million years ago the first mammoths entered North America. These mammoths came from Eurasia, crossing the Bering Strait at a time when sea level was lower than today. The first mammoths from Eurasia belonged to a species called M. meridionalis. The descendants of this species of mammoth included both the Columbian and Jefferson's mammoths. The woolly mammoths evolved in Eurasia and came over the Bering Strait much later (perhaps less than 500,000 years ago). Approximately 11,000 years ago all species of mammoths went extinct in North America. Find out more about this extinction. Although only bone and teeth of mammoths are preserved in the Midwestern U.S., the Illinois State Museum also has a sample of mammoth hair from Siberia. This photograph shows a sample of hair from the Yuribei Mammoth. The Yuribei Mammoth was found along the Yuribei River on the Gyda Penisula, NW Siberia, Russia. It is a young adult, female mammoth. She was covered with a long and thick, brownish haircoat, an example of which is shown above. Radiocarbon dating indicates that she lived about 11,000 years ago. The Yuribei Mammoth was collected in 1979 by a multidisciplinary team representing three institutes of the USSR Academy of Sciences. The hair sample shown here was presented to the Illinois State Museum in 1991 by Dr. Gennady Baryshnikov, of the History of Faunas Department of the Zoological Insitute, Russian Academy of Sciences, St. Petersburg.

http://exhibits.museum.state.il.us/exhibits/larson/mammuthus.html

Return to Mathematics Index Pitra, Barbara Marconi Community Academy The children will name, identify and categorize the geometric shapes (triangle, square and rectangle) by the number and length of the The students will identify triangles, squares and rectangles in their Through observation, comparison and manipulation the children will construct pictures, shapes and patterns with the triangle and square. Using game formats the students will gain practice in recognizing and naming geometric shapes. 1. various sized plexiglas triangles, squares and rectangles with magnetic tape applied to the backs of the shapes so they will stick to the blackboard 2. 20 triangles and squares made of plexiglas or paper (2" on a side) 3. a ditto of triangles and squares for each child (2" on a side) 6. overhead projector 7. "Color and Shape Bingo" by Trend Co. 8. 1 deck "I have ....who has...." cards This presentation is appropriate for use with primary children. Using the plexiglas shapes the students will discuss similarities and differences as well as various ways in which the objects could be categorized. Elicit from the children that a triangle has 3 sides, a square has 4 equal sides and a rectangle has 2 long sides and 2 short sides. On the blackboard make a category heading for each shape and have several students go to the board and place the magnetic shapes under the correct heading. Review names and characteristics of each Why did you place that object there? What are the characteristics of that shape? What is your shape called? Ask the children to name these geometric shapes as they identify them in the classroom and their environment. The children will cut out squares and triangles from ditto sheets in several different colors. All of the new shapes and patterns that will be constructed will be shown on the overhead projector. This will help the child who has difficulty seeing patterns or because of poor eye-hand coordination. Shapes must not overlap or cover other shapes. They must line up evenly. Ask the students to show that: squares can make bigger squares, squares can make rectangles, triangles can make bigger triangles, triangles can make hexagons but our triangles cannot make squares or rectangles (we didn't have right angles). Allow the students to experiment with color and shape in making new figures, shapes and patterns. Demonstrate some of the creations on the overhead. After the class has had adequate time to experiment, request that they glue their "favorite" on to a piece of Play Trend Co.'s "Color and Shape Bingo" or make a bingo game. The "I have ... who has ..." card game is teacher made. It consists of a deck of cards, one for each student. The cards may be used to drill many topics. For example, the child reads his card, "I have a red square, who has a blue rectangle." The child who has the blue rectangle then reads his card, etc. The last card read is the winner.

http://mypages.iit.edu/~smile/ma8714.html

French: Chicken Pox Connect to Your Teaching Reflect on Your Practice As you reflect on these questions, write down your responses or discuss them as a group. - How might you teach a new language to young students with limited literacy in their first language? How is this different from and similar to working with students who are already literate in one language? - How might you use children's or young adult literature with your students? - What immersion teaching strategies can be used in other elementary, middle, and high school models? - In classroom interactions, how do you balance the use of the target language and the use of English for students who have a limited vocabulary but are eager to communicate? Watch Other Videos Watch other videos in the Teaching Foreign Languages K-12 library for more examples of teaching methodologies like those you've just seen. Note: All videos in this series are subtitled in English. Communicating About Sports (Chinese) illustrates oral and written (character) language recognition of new vocabulary, and Holidays and Seasons (German) shows the integration of songs into a lesson. Put It Into Practice Try these ideas in your classroom. - Introduce children's literature to help students understand, reproduce, and recall language in context. A story's narrative -- beginning, conflict/problem, developments toward a resolution, and conclusion -- can help advance the meaning of new vocabulary. Children's books usually have illustrations that help readers make sense of unfamiliar words. Mr. Scott's selection, Arthur a la varicelle, appeals to children because they can relate to the character and to the childhood illness. Mr. Scott did frequent comprehension checks and allowed students to "take over" some of the story with their predictions and solutions. You could also use children's stories with some middle and high school students; success depends upon the dynamics of the group and their interest in the story itself. Students with greater language proficiency can also use the illustrations to lead the story reading with classmates. You can devote time on an ongoing basis to reading sections of longer stories. Many popular children's books are available in numerous languages. - Incorporate songs into lessons to reinforce and introduce authentic language, choosing songs that are appropriate for the grade level and the topics you teach. "La varicelle" ("The chicken pox") mentions many parts of the body and repeats phrases with unusual words -- such as itchy, scratch, and jiggle -- that quickly become familiar. These words are fun to say in French, are quickly internalized, and, when spoken using facial and hand gestures, provide an opportunity for kinesthetic and rhythmic learning. Mr. Scott taught the song by first playing it and acting out the meaning himself, and then inviting students to imitate (or reproduce) his gestures and sing along. He also showed the text to students, an activity that could be done sooner with older students. Songs can be used with all age groups, although some classes may be resistant to singing at first. In those cases you may wish to concentrate on the lyrics and not require students to sing at all. Students can act out the lyrics with gestures or movement instead.

http://www.learner.org/libraries/tfl/french/scott/connect.html

Two groups of researchers, led by Pennsylvania State University (USA), were first observed what happens in the first moments in which a black hole absorbs a star. What is surprising about this finding is that it provides a unique opportunity to study how bright the relativistic jets of matter that is issued at the beginning of the phenomenon. “Until now, this is a unique event. Although it has long been expected that such events should occur, the glow it emits is a surprise, “said Jamie A. SINC Kennea, a researcher at Pennsylvania State University and coauthor of the study published in the latest issue of the journal Nature . Scientists have determined that the black hole is at the center of a galaxy at a distance such that the light of this phenomenon took about 4 billion years to reach. Black holes are common in the centers of galaxies. The Milky Way is home to one of about 2 million times the mass of our sun Its powerful gravitational fields create strong gradients that can destroy stars that go to several million miles of it and produce a flash of ultraviolet light and X rays . “This is what we believe happened to the star absorbed in this case. The result of this process may have been observed several times, but until now had never seen beginning, “said Kennea. A jet of particles ultrafast What they have found the research teams is that the accretion-growth by adding materials to the star that has been affected by the absorption of the black hole, relativistic jet has occurred, a result not predicted by the previous theoretical models observation. “As scientists we refer to relativistic jets, means that the particles of the jet of matter moving near the speed of light. For these speeds is necessary to use the theory of relativity Einstein (hence the term ‘relativistic’) to describe the physics of the reaction.Classical or Newtonian physics does not work at these speeds. That is, the material in the jet moving very fast, about a billion miles per hour, “Kennea. DN Burrows et al. “Relativistic jet activity from the tidal disruption of a star by a massive black hole”, Nature 476: 421-424, August 25, 2011. doi: 10.1038/nature10374. |Category: Astronomy and Astrophysics||Tags: black hole|

http://www.scienceknowledge.org/2011/08/26/the-brightness-of-a-black-hole-to-absorb-a-surprise-star-scientists/

Mathematics » High School: Geometry » Similarity, Right Triangles, & Trigonometry Standards in this domain: Understand similarity in terms of similarity transformations - CCSS.Math.Content.HSG-SRT.A.1 Verify experimentally the properties of dilations given by a center and a scale factor: - CCSS.Math.Content.HSG-SRT.A.1a A dilation takes a line not passing through the center of the dilation to a parallel line, and leaves a line passing through the center unchanged. - CCSS.Math.Content.HSG-SRT.A.1b The dilation of a line segment is longer or shorter in the ratio given by the scale factor. - CCSS.Math.Content.HSG-SRT.A.2 Given two figures, use the definition of similarity in terms of similarity transformations to decide if they are similar; explain using similarity transformations the meaning of similarity for triangles as the equality of all corresponding pairs of angles and the proportionality of all corresponding pairs of sides. - CCSS.Math.Content.HSG-SRT.A.3 Use the properties of similarity transformations to establish the AA criterion for two triangles to be similar. Prove theorems involving similarity - CCSS.Math.Content.HSG-SRT.B.4 Prove theorems about triangles. Theorems include: a line parallel to one side of a triangle divides the other two proportionally, and conversely; the Pythagorean Theorem proved using triangle similarity. - CCSS.Math.Content.HSG-SRT.B.5 Use congruence and similarity criteria for triangles to solve problems and to prove relationships in geometric figures. Define trigonometric ratios and solve problems involving right triangles - CCSS.Math.Content.HSG-SRT.C.6 Understand that by similarity, side ratios in right triangles are properties of the angles in the triangle, leading to definitions of trigonometric ratios for acute angles. - CCSS.Math.Content.HSG-SRT.C.7 Explain and use the relationship between the sine and cosine of complementary angles. - CCSS.Math.Content.HSG-SRT.C.8 Use trigonometric ratios and the Pythagorean Theorem to solve right triangles in applied problems.★ Apply trigonometry to general triangles - CCSS.Math.Content.HSG-SRT.D.9 (+) Derive the formula A = 1/2 ab sin(C) for the area of a triangle by drawing an auxiliary line from a vertex perpendicular to the opposite side. - CCSS.Math.Content.HSG-SRT.D.10 (+) Prove the Laws of Sines and Cosines and use them to solve problems. - CCSS.Math.Content.HSG-SRT.D.11 (+) Understand and apply the Law of Sines and the Law of Cosines to find unknown measurements in right and non-right triangles (e.g., surveying problems, resultant forces).

http://www.corestandards.org/Math/Content/HSG/SRT

This article outlines some of the benefits of using dice games in the classroom, especially as a tool for formative assessment. Bernard's article reminds us of the richness of using dice for number, shape and probability. Simple dice and spinners tool for experiments. A game has a special dice with a colour spot on each face. These three pictures show different views of the same dice. What colour is opposite blue? A maths-based Football World Cup simulation for teachers and students to use. Use the interactivity or play this dice game yourself. How could you make it fair? In this game you throw two dice and find their total, then move the appropriate counter to the right. Which counter reaches the purple box first? Is this what you would expect? There are nasty versions of this dice game but we'll start with the nice ones... Who said that adding, subtracting, multiplying and dividing couldn't be fun? Investigate the sum of the numbers on the top and bottom faces of a line of three dice. What do you notice? In this game, you can add, subtract, multiply or divide the numbers on the dice. Which will you do so that you get to the end of the number line first? Dotty Six is a simple dice game that you can adapt in many ways. Roll two red dice and a green dice. Add the two numbers on the red dice and take away the number on the green. What are all the different possibilities that could come up? If you had any number of ordinary dice, what are the possible ways of making their totals 6? What would the product of the dice be Can you see why 2 by 2 could be 5? Can you predict what 2 by 10 Have a go at this game which involves throwing two dice and adding their totals. Where should you place your counters to be more likely to win? The idea of this game is to add or subtract the two numbers on the dice and cover the result on the grid, trying to get a line of three. Are there some numbers that are good to aim for? A game for 1 person. Can you work out how the dice must be rolled from the start position to the finish? Play on line. When you throw two regular, six-faced dice you have more chance of getting one particular result than any other. What result would that be? Why is this? I throw three dice and get 5, 3 and 2. Add the scores on the three dice. What do you get? Now multiply the scores. What do you notice?

http://nrich.maths.org/5896/index

U.S. District Court Western District of Missouri History of Jury DutyHistory of the Jury | History of the Grand Jury By the time the United States Constitution and the Bill of Rights were drafted and ratified, the institution of trial by jury was almost universally revered, so revered that its history had been traced back to Magna Carta. The jury began in the form of a grand or presentment jury with the role of inquest and was started by Frankish conquerors to discover the King's rights. Henry II regularized this type of proceeding to establish royal control over the machinery of justice, first in civil trials and then in criminal trials. Trial by petit jury was not employed at least until the reign of Henry III, in which the jury was first essentially a body of witnesses, called for their knowledge of the case; not until the reign of Henry VI did it become the trier of evidence. It was during the Seventeenth Century that the jury emerged as a safeguard for the criminally accused. Thus, in the Eighteenth Century, Blackstone could commemorate the institution as part of a ''strong and two-fold barrier . . . between the liberties of the people and the prerogative of the crown'' because ''the truth of every accusation . . . . [must] be confirmed by the unanimous suffrage of twelve of his equals and neighbors indifferently chosen and superior to all suspicion.'' The right was guaranteed in the constitutions of the original 13 States, was guaranteed in the body of the Constitution and in the Sixth Amendment, and the constitution of every State entering the Union thereafter in one form or another protected the right to jury trial in criminal cases. ''Those who emigrated to this country from England brought with them this great privilege 'as their birthright and inheritance, as a part of that admirable common law which had fenced around and interposed barriers on every side against the approaches of arbitrary power.''' ''The guarantees of jury trial in the Federal and State Constitutions reflect a profound judgment about the way in which law should be enforced and justice administered. A right to jury trial is granted to criminal defendants in order to prevent oppression by the Government. Those who wrote our constitutions knew from history and experience that it was necessary to protect against unfounded criminal charges brought to eliminate enemies and against judges too responsive to the voice of higher authority. The framers of the constitutions strove to create an independent judiciary but insisted upon further protection against arbitrary action. Providing an accused with the right to be tried by a jury of his peers gave him an inestimable safeguard against the corrupt overzealous prosecutor and against the compliant, biased, or eccentric judge. . . . [T]he jury trial provisions . . . reflect a fundamental decision about the exercise of official power--a reluctance to entrust plenary powers over the life and liberty of the citizen to one judge or to a group of judges. Fear of unchecked power . . . found expression in the criminal law in this insistence upon community participation in the determination of guilt or innocence.'' A Grand Jury derives its name from the fact that it usually has a greater number of jurors than a trial (petit) jury. One of the earliest concepts of Grand Juries dates back to early Greece where the Athenians used an accusatory body. In early Briton, the Saxons also used something similar to a Grand Jury System. During the years 978 to 1016, one of the Dooms (laws) stated that for each 100 men, 12 were to be named to act as an accusing body. They were cautioned "not to accuse an innocent man or spare a guilty one." The Grand Jury can also be traced to the time of the Norman conquest of England in 1066. There is evidence that the courts of that time summoned a body of sworn neighbors to present crimes that had come to their knowledge. Since the members of that accusing jury were selected from small jurisdictions, it was natural that they could present accusations based on their personal knowledge. Historians agree that the Assize [court session or assembly] of Clarendon in 1166 provided the ground work for our present Grand Jury system. During the reign of Henry II (1154-1189), to regain for the crown the powers usurped by Thomas Becket, Chancellor of England, 12 "good and lawful men" in each village were assembled to reveal the names of those suspected of crimes. It was during this same period that juries were divided into two types, civil and criminal, with the development of each influencing the other. The oath taken by these jurors provided that they would carry out their duties faithfully, that they would aggrieve no one through enmity nor deference to anyone through love, and that they would conceal those things that they had heard. By the year 1290, these accusing juries were given the authority to inquire into the maintenance of bridges and highways, defects of jails, and whether the Sheriff had kept in jail anyone who should have been brought before the justices. "Le Grand Inquest" evolved during the reign of Edward III (1368), when the "accusatory jury" was increased in number from 12 to 23, with a majority vote necessary to indict anyone accused of crime. In America, the Massachusetts Bay Colony impaneled the first Grand Jury in 1635 to consider cases of murder, robbery and wife beating. As early as 1700, the value of the Grand Jury was recognized as opposing the Royalists. These colonial Grand Juries expressed their independence by refusing to indict leaders of the Stamp Act (1765), and refusing to bring libel charges against the editors of the Boston Gazette (1765). A union with other colonies to oppose British taxes was supported by the Philadelphia Grand Jury in 1770. By the end of the Colonial Period, the Grand Jury had become an indispensable adjunct of Government: "they proposed new laws, protested against abuses in government, and wielded the tremendous authority in their power to determine who should and who should not face trial."

http://www.mow.uscourts.gov/district/jury/jury_history.html

“The Greenland Ice Sheet could contribute up to seven meters of global sea-level rise if it were to melt,” says OSU marine geologist Joseph Stoner. “We don’t know if it’s going to melt, but that’s how much water is in the ice sheet. Therefore, we need to better understand the processes at work.” In search of that understanding, Stoner and researchers at the University of Wisconsin-Madison are studying sediments flowing seaward in streams and rivers on the island’s southern tip. Those sediments—remnants of bedrock pulverized over eons by grinding glaciers and rushing rivers—hold clues to the ice sheet’s history across geologic time. Scientists know that the 680,000-cubic mile chunk of snow, compressed from white to crystalline blue over many millennia, is receding. Satellite images from the past several decades show significant shrinkage. What isn’t known is the speed of melting or the extent that melting might take in coming years. By studying Greenland’s past with support from the National Science Foundation, Stoner and his colleagues hope to bring its future into clearer focus. “The key to understanding the Greenland Ice Sheet is to use the natural record of past variability as a sort of a manual as to what it could do in the future,” says Stoner. “We’re trying to use the natural geological archive to test how the ice sheet works.” To recreate the ice sheet’s prehistoric behavior, he and his graduatestudents collected sediment samples. Tracing the origins of these silts and sands should tell the researchers where the island was exposed during “interglacial” periods—warm stretches between ice ages—and where it lay buried beneath tons of frozen snow during colder periods. The markers that will reveal these ancient patterns are both chemical and magnetic. To read magnetic profiles of sediments, Stoner’s lab recently acquired a new-generation instrument: a superconducting magnetometer for measuring the magnetic properties and composition of rocks. Instead of using liquid helium as a coolant like old-style cryogenic magnetometers do, this one compresses helium gas till it reaches 3.5 degrees Kelvin, “Just a little above absolute zero,” Stoner says. “It works through superconductivity, which only happens at extremely cold temperatures.” Stoner’s findings could cause scientists to rethink Greenland’s role in climate-change scenarios. Article courtesy of the College of Earth, Ocean, and Atmospheric Sciences 2010 Research Highlights Photo courtesy of Terra Magazine

http://oregonstate.edu/marinesciences/greenland-ice-sheet