Update models

README.md

@@ -0,0 +1,82 @@
+---
+license: other
+tags:
+- stable-diffusion
+- text-to-image
+- core-ml
+---
+
+# Stable Diffusion XL v0.9 Model Card
+
+This model was generated using [Apple’s repository](https://github.com/apple/ml-stable-diffusion) which has [ASCL](https://github.com/apple/ml-stable-diffusion/blob/main/LICENSE.md).
+
+This model card focuses on the model associated with the Stable Diffusion XL v0.9 Base model, codebase available [here](https://github.com/Stability-AI/generative-models).
+
+SDXL v0.9 consists of a two-step pipeline for latent diffusion: 
+First, we use a base model to generate latents of the desired output size. 
+In the second step, we use a specialized high-resolution model and apply a technique called SDEdit (https://arxiv.org/abs/2108.01073, also known as "img2img") 
+to the latents generated in the first step, using the same prompt.
+
+Only the base model is included here.
+
+These weights here have been converted to Core ML for use on Apple Silicon hardware.
+
+There are 2 variants of the Core ML weights:
+
+```
+coreml-stable-diffusion-xl-v0-9-base
+└── original
+ ├── compiled # Swift inference, "original" attention
+ └── packages # Python inference, "original" 
+```
+
+### Model Description
+
+- **Developed by:** Stability AI
+- **Model type:** Diffusion-based text-to-image generative model
+- **License:** [SDXL 0.9 Research License](https://huggingface.co/stabilityai/stable-diffusion-xl-base-0.9/blob/main/LICENSE.md)
+- **Model Description:** This is a model that can be used to generate and modify images based on text prompts. It is a [Latent Diffusion Model](https://arxiv.org/abs/2112.10752) that uses two fixed, pretrained text encoders ([OpenCLIP-ViT/G](https://github.com/mlfoundations/open_clip) and [CLIP-ViT/L](https://github.com/openai/CLIP/tree/main)).
+- **Resources for more information:** [GitHub Repository](https://github.com/Stability-AI/generative-models) [SDXL paper on arXiv](https://arxiv.org/abs/2307.01952).
+
+### Model Sources
+
+<!-- Provide the basic links for the model. -->
+
+- **Repository:** https://github.com/Stability-AI/generative-models
+- **Demo [optional]:** https://clipdrop.co/stable-diffusion
+
+## Uses
+
+### Direct Use
+
+The model is intended for research purposes only. Possible research areas and tasks include
+
+- Generation of artworks and use in design and other artistic processes.
+- Applications in educational or creative tools.
+- Research on generative models.
+- Safe deployment of models which have the potential to generate harmful content.
+- Probing and understanding the limitations and biases of generative models.
+
+Excluded uses are described below.
+
+### Out-of-Scope Use
+
+The model was not trained to be factual or true representations of people or events, and therefore using the model to generate such content is out-of-scope for the abilities of this model.
+
+## Limitations and Bias
+
+### Limitations
+
+- The model does not achieve perfect photorealism
+- The model cannot render legible text
+- The model struggles with more difficult tasks which involve compositionality, such as rendering an image corresponding to “A red cube on top of a blue sphere”
+- Faces and people in general may not be generated properly.
+- The autoencoding part of the model is lossy.
+
+### Bias
+While the capabilities of image generation models are impressive, they can also reinforce or exacerbate social biases.
+
+## Evaluation
+![comparison](https://huggingface.co/stabilityai/stable-diffusion-xl-base-0.9/resolve/main/comparison.png)
+The chart above evaluates user preference for SDXL (with and without refinement) over Stable Diffusion 1.5 and 2.1. 
+The SDXL base model performs significantly better than the previous variants, and the model combined with the refinement module achieves the best overall performance.

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