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-# Generative Models by Stability AI
-
-![sample1](assets/000.jpg)
-
-## News
-
-**November 21, 2023**
-
-- We are releasing Stable Video Diffusion, an image-to-video model, for research purposes:
-    - [SVD](https://huggingface.co/stabilityai/stable-video-diffusion-img2vid): This model was trained to generate 14
-      frames at resolution 576x1024 given a context frame of the same size. 
-      We use the standard image encoder from SD 2.1, but replace the decoder with a temporally-aware `deflickering decoder`.
-    - [SVD-XT](https://huggingface.co/stabilityai/stable-video-diffusion-img2vid-xt): Same architecture as `SVD` but finetuned
-      for 25 frame generation.
-    - We provide a streamlit demo `scripts/demo/video_sampling.py` and a standalone python script `scripts/sampling/simple_video_sample.py` for inference of both models.
-    - Alongside the model, we release a [technical report](https://stability.ai/research/stable-video-diffusion-scaling-latent-video-diffusion-models-to-large-datasets).
-
-  ![tile](assets/tile.gif)
-
-**July 26, 2023**
-
-- We are releasing two new open models with a
-  permissive [`CreativeML Open RAIL++-M` license](model_licenses/LICENSE-SDXL1.0) (see [Inference](#inference) for file
-  hashes):
-    - [SDXL-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0): An improved version
-      over `SDXL-base-0.9`.
-    - [SDXL-refiner-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0): An improved version
-      over `SDXL-refiner-0.9`.
-
-![sample2](assets/001_with_eval.png)
-
-**July 4, 2023**
-
-- A technical report on SDXL is now available [here](https://arxiv.org/abs/2307.01952).
-
-**June 22, 2023**
-
-- We are releasing two new diffusion models for research purposes:
-    - `SDXL-base-0.9`: The base model was trained on a variety of aspect ratios on images with resolution 1024^2. The
-      base model uses [OpenCLIP-ViT/G](https://github.com/mlfoundations/open_clip)
-      and [CLIP-ViT/L](https://github.com/openai/CLIP/tree/main) for text encoding whereas the refiner model only uses
-      the OpenCLIP model.
-    - `SDXL-refiner-0.9`: The refiner has been trained to denoise small noise levels of high quality data and as such is
-      not expected to work as a text-to-image model; instead, it should only be used as an image-to-image model.
-
-If you would like to access these models for your research, please apply using one of the following links:
-[SDXL-0.9-Base model](https://huggingface.co/stabilityai/stable-diffusion-xl-base-0.9),
-and [SDXL-0.9-Refiner](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-0.9).
-This means that you can apply for any of the two links - and if you are granted - you can access both.
-Please log in to your Hugging Face Account with your organization email to request access.
-**We plan to do a full release soon (July).**
-
-## The codebase
-
-### General Philosophy
-
-Modularity is king. This repo implements a config-driven approach where we build and combine submodules by
-calling `instantiate_from_config()` on objects defined in yaml configs. See `configs/` for many examples.
-
-### Changelog from the old `ldm` codebase
-
-For training, we use [PyTorch Lightning](https://lightning.ai/docs/pytorch/stable/), but it should be easy to use other
-training wrappers around the base modules. The core diffusion model class (formerly `LatentDiffusion`,
-now `DiffusionEngine`) has been cleaned up:
-
-- No more extensive subclassing! We now handle all types of conditioning inputs (vectors, sequences and spatial
-  conditionings, and all combinations thereof) in a single class: `GeneralConditioner`,
-  see `sgm/modules/encoders/modules.py`.
-- We separate guiders (such as classifier-free guidance, see `sgm/modules/diffusionmodules/guiders.py`) from the
-  samplers (`sgm/modules/diffusionmodules/sampling.py`), and the samplers are independent of the model.
-- We adopt the ["denoiser framework"](https://arxiv.org/abs/2206.00364) for both training and inference (most notable
-  change is probably now the option to train continuous time models):
-    * Discrete times models (denoisers) are simply a special case of continuous time models (denoisers);
-      see `sgm/modules/diffusionmodules/denoiser.py`.
-    * The following features are now independent: weighting of the diffusion loss
-      function (`sgm/modules/diffusionmodules/denoiser_weighting.py`), preconditioning of the
-      network (`sgm/modules/diffusionmodules/denoiser_scaling.py`), and sampling of noise levels during
-      training (`sgm/modules/diffusionmodules/sigma_sampling.py`).
-- Autoencoding models have also been cleaned up.
-
-## Installation:
-
-<a name="installation"></a>
-
-#### 1. Clone the repo
-
-```shell
-git clone [email protected]:Stability-AI/generative-models.git
-cd generative-models
-```
-
-#### 2. Setting up the virtualenv
-
-This is assuming you have navigated to the `generative-models` root after cloning it.
-
-**NOTE:** This is tested under `python3.10`. For other python versions, you might encounter version conflicts.
-
-**PyTorch 2.0**
-
-```shell
-# install required packages from pypi
-python3 -m venv .pt2
-source .pt2/bin/activate
-pip3 install -r requirements/pt2.txt
-```
-
-#### 3. Install `sgm`
-
-```shell
-pip3 install .
-```
-
-#### 4. Install `sdata` for training
-
-```shell
-pip3 install -e git+https://github.com/Stability-AI/datapipelines.git@main#egg=sdata
-```
-
-## Packaging
-
-This repository uses PEP 517 compliant packaging using [Hatch](https://hatch.pypa.io/latest/).
-
-To build a distributable wheel, install `hatch` and run `hatch build`
-(specifying `-t wheel` will skip building a sdist, which is not necessary).
-
-```
-pip install hatch
-hatch build -t wheel
-```
-
-You will find the built package in `dist/`. You can install the wheel with `pip install dist/*.whl`.
-
-Note that the package does **not** currently specify dependencies; you will need to install the required packages,
-depending on your use case and PyTorch version, manually.
-
-## Inference
-
-We provide a [streamlit](https://streamlit.io/) demo for text-to-image and image-to-image sampling
-in `scripts/demo/sampling.py`.
-We provide file hashes for the complete file as well as for only the saved tensors in the file (
-see [Model Spec](https://github.com/Stability-AI/ModelSpec) for a script to evaluate that).
-The following models are currently supported:
-
-- [SDXL-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0)
-  ```
-  File Hash (sha256): 31e35c80fc4829d14f90153f4c74cd59c90b779f6afe05a74cd6120b893f7e5b
-  Tensordata Hash (sha256): 0xd7a9105a900fd52748f20725fe52fe52b507fd36bee4fc107b1550a26e6ee1d7
-  ```
-- [SDXL-refiner-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0)
-  ```
-  File Hash (sha256): 7440042bbdc8a24813002c09b6b69b64dc90fded4472613437b7f55f9b7d9c5f
-  Tensordata Hash (sha256): 0x1a77d21bebc4b4de78c474a90cb74dc0d2217caf4061971dbfa75ad406b75d81
-  ```
-- [SDXL-base-0.9](https://huggingface.co/stabilityai/stable-diffusion-xl-base-0.9)
-- [SDXL-refiner-0.9](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-0.9)
-- [SD-2.1-512](https://huggingface.co/stabilityai/stable-diffusion-2-1-base/blob/main/v2-1_512-ema-pruned.safetensors)
-- [SD-2.1-768](https://huggingface.co/stabilityai/stable-diffusion-2-1/blob/main/v2-1_768-ema-pruned.safetensors)
-
-**Weights for SDXL**:
-
-**SDXL-1.0:**
-The weights of SDXL-1.0 are available (subject to
-a [`CreativeML Open RAIL++-M` license](model_licenses/LICENSE-SDXL1.0)) here:
-
-- base model: https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0/
-- refiner model: https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-1.0/
-
-**SDXL-0.9:**
-The weights of SDXL-0.9 are available and subject to a [research license](model_licenses/LICENSE-SDXL0.9).
-If you would like to access these models for your research, please apply using one of the following links:
-[SDXL-base-0.9 model](https://huggingface.co/stabilityai/stable-diffusion-xl-base-0.9),
-and [SDXL-refiner-0.9](https://huggingface.co/stabilityai/stable-diffusion-xl-refiner-0.9).
-This means that you can apply for any of the two links - and if you are granted - you can access both.
-Please log in to your Hugging Face Account with your organization email to request access.
-
-After obtaining the weights, place them into `checkpoints/`.
-Next, start the demo using
-
-```
-streamlit run scripts/demo/sampling.py --server.port <your_port>
-```
-
-### Invisible Watermark Detection
-
-Images generated with our code use the
-[invisible-watermark](https://github.com/ShieldMnt/invisible-watermark/)
-library to embed an invisible watermark into the model output. We also provide
-a script to easily detect that watermark. Please note that this watermark is
-not the same as in previous Stable Diffusion 1.x/2.x versions.
-
-To run the script you need to either have a working installation as above or
-try an _experimental_ import using only a minimal amount of packages:
-
-```bash
-python -m venv .detect
-source .detect/bin/activate
-
-pip install "numpy>=1.17" "PyWavelets>=1.1.1" "opencv-python>=4.1.0.25"
-pip install --no-deps invisible-watermark
-```
-
-To run the script you need to have a working installation as above. The script
-is then useable in the following ways (don't forget to activate your
-virtual environment beforehand, e.g. `source .pt1/bin/activate`):
-
-```bash
-# test a single file
-python scripts/demo/detect.py <your filename here>
-# test multiple files at once
-python scripts/demo/detect.py <filename 1> <filename 2> ... <filename n>
-# test all files in a specific folder
-python scripts/demo/detect.py <your folder name here>/*
-```
-
-## Training:
-
-We are providing example training configs in `configs/example_training`. To launch a training, run
-
-```
-python main.py --base configs/<config1.yaml> configs/<config2.yaml>
-```
-
-where configs are merged from left to right (later configs overwrite the same values).
-This can be used to combine model, training and data configs. However, all of them can also be
-defined in a single config. For example, to run a class-conditional pixel-based diffusion model training on MNIST,
-run
-
-```bash
-python main.py --base configs/example_training/toy/mnist_cond.yaml
-```
-
-**NOTE 1:** Using the non-toy-dataset
-configs `configs/example_training/imagenet-f8_cond.yaml`, `configs/example_training/txt2img-clipl.yaml`
-and `configs/example_training/txt2img-clipl-legacy-ucg-training.yaml` for training will require edits depending on the
-used dataset (which is expected to stored in tar-file in
-the [webdataset-format](https://github.com/webdataset/webdataset)). To find the parts which have to be adapted, search
-for comments containing `USER:` in the respective config.
-
-**NOTE 2:** This repository supports both `pytorch1.13` and `pytorch2`for training generative models. However for
-autoencoder training as e.g. in `configs/example_training/autoencoder/kl-f4/imagenet-attnfree-logvar.yaml`,
-only `pytorch1.13` is supported.
-
-**NOTE 3:** Training latent generative models (as e.g. in `configs/example_training/imagenet-f8_cond.yaml`) requires
-retrieving the checkpoint from [Hugging Face](https://huggingface.co/stabilityai/sdxl-vae/tree/main) and replacing
-the `CKPT_PATH` placeholder in [this line](configs/example_training/imagenet-f8_cond.yaml#81). The same is to be done
-for the provided text-to-image configs.
-
-### Building New Diffusion Models
-
-#### Conditioner
-
-The `GeneralConditioner` is configured through the `conditioner_config`. Its only attribute is `emb_models`, a list of
-different embedders (all inherited from `AbstractEmbModel`) that are used to condition the generative model.
-All embedders should define whether or not they are trainable (`is_trainable`, default `False`), a classifier-free
-guidance dropout rate is used (`ucg_rate`, default `0`), and an input key (`input_key`), for example, `txt` for
-text-conditioning or `cls` for class-conditioning.
-When computing conditionings, the embedder will get `batch[input_key]` as input.
-We currently support two to four dimensional conditionings and conditionings of different embedders are concatenated
-appropriately.
-Note that the order of the embedders in the `conditioner_config` is important.
-
-#### Network
-
-The neural network is set through the `network_config`. This used to be called `unet_config`, which is not general
-enough as we plan to experiment with transformer-based diffusion backbones.
-
-#### Loss
-
-The loss is configured through `loss_config`. For standard diffusion model training, you will have to
-set `sigma_sampler_config`.
-
-#### Sampler config
-
-As discussed above, the sampler is independent of the model. In the `sampler_config`, we set the type of numerical
-solver, number of steps, type of discretization, as well as, for example, guidance wrappers for classifier-free
-guidance.
-
-### Dataset Handling
-
-For large scale training we recommend using the data pipelines from
-our [data pipelines](https://github.com/Stability-AI/datapipelines) project. The project is contained in the requirement
-and automatically included when following the steps from the [Installation section](#installation).
-Small map-style datasets should be defined here in the repository (e.g., MNIST, CIFAR-10, ...), and return a dict of
-data keys/values,
-e.g.,
-
-```python
-example = {"jpg": x,  # this is a tensor -1...1 chw
-           "txt": "a beautiful image"}
-```
-
-where we expect images in -1...1, channel-first format.
+---
+title: Stable Video Diffusion
+emoji: 📺
+colorFrom: purple
+colorTo: purple
+sdk: gradio
+sdk_version: 4.4.0
+app_file: app.py
+pinned: false
+license: other
+---