---
license: creativeml-openrail-m
language: en
base_model: stable-diffusion-xl
tags:
  - stable-diffusion
  - text-to-image
  - generative
  - diffusion
metrics:
  - FID
pipeline_tag: text-to-image
library_name: diffusers
---

# Model Card: **JuggernautXL_v8Rundiffusion**

## Overview
**JuggernautXL_v8Rundiffusion** is a cutting-edge model built for **stable diffusion-based generative tasks**. This model is optimized for creating high-quality, detailed, and cinematic images using state-of-the-art diffusion techniques. It supports multiple LoRA (Low-Rank Adaptation) integrations and excels at generating diverse styles, including realistic, artistic, and abstract images.

---

## Features
- **Model Type:** Stable Diffusion XL (SDXL)
- **Version:** v8 Rundiffusion
- **Key Capabilities:**
  - Supports intricate artistic styles, cinematic imagery, and vibrant compositions.
  - Works seamlessly with additional LoRA files for fine-tuned styles.
  - Efficient memory usage and adaptability to various VRAM configurations.
- **Applications:**
  - **Creative Design:** Digital art, concept design, and visual effects.
  - **Marketing:** Ad campaigns, social media content creation.
  - **Educational:** Demonstrations of diffusion-based generative AI.

---

## Usage
This model requires a GPU with at least **16 GB VRAM** for optimal performance.

### Example Code:
```python
from diffusers import StableDiffusionPipeline

# Load the model from Hugging Face
pipeline = StableDiffusionPipeline.from_pretrained("oieieio/juggernautXL_v8Rundiffusion")
pipeline.to("cuda")

# Generate an image
prompt = "a futuristic cityscape, cinematic lighting, ultra-detailed"
image = pipeline(prompt).images[0]

# Save the image
image.save("output.png")